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The classification of finite simple groups (CFSG), first announced in 1983 but only fully completed in 2004, is one of the monumental achievements of twentieth century mathematics. Spanning hundreds of papers and tens of thousands of pages, it has been called the “enormous theorem”. A “second generation” proof of the theorem is nearly completed which is a little shorter (estimated at about five thousand pages in length), but currently there is no reasonably sized proof of the classification.

An important precursor of the CFSG is the Feit-Thompson theorem from 1962-1963, which asserts that every finite group of odd order is solvable, or equivalently that every non-abelian finite simple group has even order. This is an immediate consequence of CFSG, and conversely the Feit-Thompson theorem is an essential starting point in the proof of the classification, since it allows one to reduce matters to groups of even order for which key additional tools (such as the Brauer-Fowler theorem) become available. The original proof of the Feit-Thompson theorem is 255 pages long, which is significantly shorter than the proof of the CFSG, but still far from short. While parts of the proof of the Feit-Thompson theorem have been simplified (and it has recently been converted, after six years of effort, into an argument that has been verified by the proof assistant Coq), the available proofs of this theorem are still extremely lengthy by any reasonable standard.

However, there is a significantly simpler special case of the Feit-Thompson theorem that was established previously by Suzuki in 1957, which was influential in the proof of the more general Feit-Thompson theorem (and thus indirectly to the proof of CFSG). Define a CA-group to be a group ${G}$ with the property that the centraliser ${C_G(x) := \{ g \in G: gx=xg \}}$ of any non-identity element ${x \in G}$ is abelian; equivalently, the commuting relation ${x \sim y}$ (defined as the relation that holds when ${x}$ commutes with ${y}$, thus ${xy=yx}$) is an equivalence relation on the non-identity elements ${G \backslash \{1\}}$ of ${G}$. Trivially, every abelian group is CA. A non-abelian example of a CA-group is the ${ax+b}$ group of invertible affine transformations ${x \mapsto ax+b}$ on a field ${F}$. A little less obviously, the special linear group ${SL_2(F_q)}$ over a finite field ${F_q}$ is a CA-group when ${q}$ is a power of two. The finite simple groups of Lie type are not, in general, CA-groups, but when the rank is bounded they tend to behave as if they were “almost CA”; the centraliser of a generic element in ${SL_d(F_q)}$, for instance, when ${d}$ is bounded and ${q}$ is large), is typically a maximal torus (because most elements in ${SL_d(F_q)}$ are regular semisimple) which is certainly abelian. In view of the CFSG, we thus see that CA or nearly CA groups form an important subclass of the simple groups, and it is thus of interest to study them separately. To this end, we have

Theorem 1 (Suzuki’s theorem on CA-groups) Every finite CA-group of odd order is solvable.

Of course, this theorem is superceded by the more general Feit-Thompson theorem, but Suzuki’s proof is substantially shorter (the original proof is nine pages) and will be given in this post. (See this survey of Solomon for some discussion of the link between Suzuki’s argument and the Feit-Thompson argument.) Suzuki’s analysis can be pushed further to give an essentially complete classification of all the finite CA-groups (of either odd or even order), but we will not pursue these matters here.

Moving even further down the ladder of simple precursors of CSFG is the following theorem of Frobenius from 1901. Define a Frobenius group to be a finite group ${G}$ which has a subgroup ${H}$ (called the Frobenius complement) with the property that all the non-trivial conjugates ${gHg^{-1}}$ of ${H}$ for ${g \in G \backslash H}$, intersect ${H}$ only at the origin. For instance the ${ax+b}$ group is also a Frobenius group (take ${H}$ to be the affine transformations that fix a specified point ${x_0 \in F}$, e.g. the origin). This example suggests that there is some overlap between the notions of a Frobenius group and a CA group. Indeed, note that if ${G}$ is a CA-group and ${H}$ is a maximal abelian subgroup of ${G}$, then any conjugate ${gHg^{-1}}$ of ${H}$ that is not identical to ${H}$ will intersect ${H}$ only at the origin (because ${H}$ and each of its conjugates consist of equivalence classes under the commuting relation ${\sim}$, together with the identity). So if a maximal abelian subgroup ${H}$ of a CA-group is its own normaliser (thus ${N(H) := \{ g \in G: gH=Hg\}}$ is equal to ${H}$), then the group is a Frobenius group.

Frobenius’ theorem places an unexpectedly strong amount of structure on a Frobenius group:

Theorem 2 (Frobenius’ theorem) Let ${G}$ be a Frobenius group with Frobenius complement ${H}$. Then there exists a normal subgroup ${K}$ of ${G}$ (called the Frobenius kernel of ${G}$) such that ${G}$ is the semi-direct product ${H \ltimes K}$ of ${H}$ and ${K}$.

Roughly speaking, this theorem indicates that all Frobenius groups “behave” like the ${ax+b}$ example (which is a quintessential example of a semi-direct product).

Note that if every CA-group of odd order was either Frobenius or abelian, then Theorem 2 would imply Theorem 1 by an induction on the order of ${G}$, since any subgroup of a CA-group is clearly again a CA-group. Indeed, the proof of Suzuki’s theorem does basically proceed by this route (Suzuki’s arguments do indeed imply that CA-groups of odd order are Frobenius or abelian, although we will not quite establish that fact here).

Frobenius’ theorem can be reformulated in the following concrete combinatorial form:

Theorem 3 (Frobenius’ theorem, equivalent version) Let ${G}$ be a group of permutations acting transitively on a finite set ${X}$, with the property that any non-identity permutation in ${G}$ fixes at most one point in ${X}$. Then the set of permutations in ${G}$ that fix no points in ${X}$, together with the identity, is closed under composition.

Again, a good example to keep in mind for this theorem is when ${G}$ is the group of affine permutations on a field ${F}$ (i.e. the ${ax+b}$ group for that field), and ${X}$ is the set of points on that field. In that case, the set of permutations in ${G}$ that do not fix any points are the non-trivial translations.

To deduce Theorem 3 from Theorem 2, one applies Theorem 2 to the stabiliser of a single point in ${X}$. Conversely, to deduce Theorem 2 from Theorem 3, set ${X := G/H = \{ gH: g \in G \}}$ to be the space of left-cosets of ${H}$, with the obvious left ${G}$-action; one easily verifies that this action is faithful, transitive, and each non-identity element ${g}$ of ${G}$ fixes at most one left-coset of ${H}$ (basically because it lies in at most one conjugate of ${H}$). If we let ${K}$ be the elements of ${G}$ that do not fix any point in ${X}$, plus the identity, then by Theorem 3 ${K}$ is closed under composition; it is also clearly closed under inverse and conjugation, and is hence a normal subgroup of ${G}$. From construction ${K}$ is the identity plus the complement of all the ${|G|/|H|}$ conjugates of ${H}$, which are all disjoint except at the identity, so by counting elements we see that

$\displaystyle |K| = |G| - \frac{|G|}{|H|}(|H|-1) = |G|/|H|.$

As ${H}$ normalises ${K}$ and is disjoint from ${K}$, we thus see that ${KH = H \ltimes K}$ is all of ${G}$, giving Theorem 2.

Despite the appealingly concrete and elementary form of Theorem 3, the only known proofs of that theorem (or equivalently, Theorem 2) in its full generality proceed via the machinery of group characters (which one can think of as a version of Fourier analysis for nonabelian groups). On the other hand, once one establishes the basic theory of these characters (reviewed below the fold), the proof of Frobenius’ theorem is very short, which gives quite a striking example of the power of character theory. The proof of Suzuki’s theorem also proceeds via character theory, and is basically a more involved version of the Frobenius argument; again, no character-free proof of Suzuki’s theorem is currently known. (The proofs of Feit-Thompson and CFSG also involve characters, but those proofs also contain many other arguments of much greater complexity than the character-based portions of the proof.)

It seems to me that the above four theorems (Frobenius, Suzuki, Feit-Thompson, and CFSG) provide a ladder of sorts (with exponentially increasing complexity at each step) to the full classification, and that any new approach to the classification might first begin by revisiting the earlier theorems on this ladder and finding new proofs of these results first (in particular, if one had a “robust” proof of Suzuki’s theorem that also gave non-trivial control on “almost CA-groups” – whatever that means – then this might lead to a new route to classifying the finite simple groups of Lie type and bounded rank). But even for the simplest two results on this ladder – Frobenius and Suzuki – it seems remarkably difficult to find any proof that is not essentially the character-based proof. (Even trying to replace character theory by its close cousin, representation theory, doesn’t seem to work unless one gives in to the temptation to take traces everywhere and put the characters back in; it seems that rather than abandon characters altogether, one needs to find some sort of “robust” generalisation of existing character-based methods.) In any case, I am recording here the standard character-based proofs of the theorems of Frobenius and Suzuki below the fold. There is nothing particularly novel here, but I wanted to collect all the relevant material in one place, largely for my own benefit.

I’ve just uploaded to the arXiv my paper “Mixing for progressions in non-abelian groups“, submitted to Forum of Mathematics, Sigma (which, along with sister publication Forum of Mathematics, Pi, has just opened up its online submission system). This paper is loosely related in subject topic to my two previous papers on polynomial expansion and on recurrence in quasirandom groups (with Vitaly Bergelson), although the methods here are rather different from those in those two papers. The starting motivation for this paper was a question posed in this foundational paper of Tim Gowers on quasirandom groups. In that paper, Gowers showed (among other things) that if ${G}$ was a quasirandom group, patterns such as ${(x,xg,xh, xgh)}$ were mixing in the sense that, for any four sets ${A,B,C,D \subset G}$, the number of such quadruples ${(x,xg,xh, xgh)}$ in ${A \times B \times C \times D}$ was equal to ${(\mu(A) \mu(B) \mu(C) \mu(D) + o(1)) |G|^3}$, where ${\mu(A) := |A| / |G|}$, and ${o(1)}$ denotes a quantity that goes to zero as the quasirandomness of the group goes to infinity. In my recent paper with Vitaly, we also considered mixing properties of some other patterns, namely ${(x,xg,gx)}$ and ${(g,x,xg,gx)}$. This paper is concerned instead with the pattern ${(x,xg,xg^2)}$, that is to say a geometric progression of length three. As observed by Gowers, by applying (a suitably quantitative version of) Roth’s theorem in (cosets of) a cyclic group, one can obtain a recurrence theorem for this pattern without much effort: if ${G}$ is an arbitrary finite group, and ${A}$ is a subset of ${G}$ with ${\mu(A) \geq \delta}$, then there are at least ${c(\delta) |G|^2}$ pairs ${(x,g) \in G}$ such that ${x, xg, xg^2 \in A}$, where ${c(\delta)>0}$ is a quantity depending only on ${\delta}$. However, this argument does not settle the question of whether there is a stronger mixing property, in that the number of pairs ${(x,g) \in G^2}$ such that ${(x,xg,xg^2) \in A \times B \times C}$ should be ${(\mu(A)\mu(B)\mu(C)+o(1)) |G|^2}$ for any ${A,B,C \subset G}$. Informally, this would assert that for ${x, g}$ chosen uniformly at random from ${G}$, the triplet ${(x, xg, xg^2)}$ should resemble a uniformly selected element of ${G^3}$ in some weak sense.

For non-quasirandom groups, such mixing properties can certainly fail. For instance, if ${G}$ is the cyclic group ${G = ({\bf Z}/N{\bf Z},+)}$ (which is abelian and thus highly non-quasirandom) with the additive group operation, and ${A = \{1,\ldots,\lfloor \delta N\rfloor\}}$ for some small but fixed ${\delta > 0}$, then ${\mu(A) = \delta + o(1)}$ in the limit ${N \rightarrow \infty}$, but the number of pairs ${(x,g) \in G^2}$ with ${x, x+g, x+2g \in A}$ is ${(\delta^2/2 + o(1)) |G|^2}$ rather than ${(\delta^3+o(1)) |G|^2}$. The problem here is that the identity ${(x+2g) = 2(x+g) - x}$ ensures that if ${x}$ and ${x+g}$ both lie in ${A}$, then ${x+2g}$ has a highly elevated likelihood of also falling in ${A}$. One can view ${A}$ as the preimage of a small ball under the one-dimensional representation ${\rho: G \rightarrow U(1)}$ defined by ${\rho(n) := e^{2\pi i n/N}}$; similar obstructions to mixing can also be constructed from other low-dimensional representations.

However, by definition, quasirandom groups do not have low-dimensional representations, and Gowers asked whether mixing for ${(x,xg,xg^2)}$ could hold for quasirandom groups. I do not know if this is the case for arbitrary quasirandom groups, but I was able to settle the question for a specific class of quasirandom groups, namely the special linear groups ${G := SL_d(F)}$ over a finite field ${F}$ in the regime where the dimension ${d}$ is bounded (but is at least two) and ${F}$ is large. Indeed, for such groups I can obtain a count of ${(\mu(A) \mu(B) \mu(C) + O( |F|^{-\min(d-1,2)/8} )) |G|^2}$ for the number of pairs ${(x,g) \in G^2}$ with ${(x, xg, xg^2) \in A \times B \times C}$. In fact, I have the somewhat stronger statement that there are ${(\mu(A) \mu(B) \mu(C) \mu(D) + O( |F|^{-\min(d-1,2)/8} )) |G|^2}$ pairs ${(x,g) \in G^2}$ with ${(x,xg,xg^2,g) \in A \times B \times C \times D}$ for any ${A,B,C,D \subset G}$.

I was also able to obtain a partial result for the length four progression ${(x,xg,xg^2, xg^3)}$ in the simpler two-dimensional case ${G = SL_2(F)}$, but I had to make the unusual restriction that the group element ${g \in G}$ was hyperbolic in the sense that it was diagonalisable over the finite field ${F}$ (as opposed to diagonalisable over the algebraic closure ${\overline{F}}$ of that field); this amounts to the discriminant of the matrix being a quadratic residue, and this holds for approximately half of the elements of ${G}$. The result is then that for any ${A,B,C,D \subset G}$, one has ${(\frac{1}{2} \mu(A) \mu(B) \mu(C) \mu(D) + o(1)) |G|^2}$ pairs ${(x,g) \in G}$ with ${g}$ hyperbolic and ${(x,xg,xg^2,xg^3) \subset A \times B \times C \times D}$. (Again, I actually show a slightly stronger statement in which ${g}$ is restricted to an arbitrary subset ${E}$ of hyperbolic elements.)

For the length three argument, the main tools used are the Cauchy-Schwarz inequality, the quasirandomness of ${G}$, and some algebraic geometry to ensure that a certain family of probability measures on ${G}$ that are defined algebraically are approximately uniformly distributed. The length four argument is significantly more difficult and relies on a rather ad hoc argument involving, among other things, expander properties related to the work of Bourgain and Gamburd, and also a “twisted” version of an argument of Gowers that is used (among other things) to establish an inverse theorem for the ${U^3}$ norm.

I give some details of these arguments below the fold.

I’ve just uploaded to the arXiv my joint paper with Vitaly Bergelson, “Multiple recurrence in quasirandom groups“, which is submitted to Geom. Func. Anal.. This paper builds upon a paper of Gowers in which he introduced the concept of a quasirandom group, and established some mixing (or recurrence) properties of such groups. A ${D}$-quasirandom group is a finite group with no non-trivial unitary representations of dimension at most ${D}$. We will informally refer to a “quasirandom group” as a ${D}$-quasirandom group with the quasirandomness parameter ${D}$ large (more formally, one can work with a sequence of ${D_n}$-quasirandom groups with ${D_n}$ going to infinity). A typical example of a quasirandom group is ${SL_2(F_p)}$ where ${p}$ is a large prime. Quasirandom groups are discussed in depth in this blog post. One of the key properties of quasirandom groups established in Gowers’ paper is the following “weak mixing” property: if ${A, B}$ are subsets of ${G}$, then for “almost all” ${g \in G}$, one has

$\displaystyle \mu( A \cap gB ) \approx \mu(A) \mu(B) \ \ \ \ \ (1)$

where ${\mu(A) := |A|/|G|}$ denotes the density of ${A}$ in ${G}$. Here, we use ${x \approx y}$ to informally represent an estimate of the form ${x=y+o(1)}$ (where ${o(1)}$ is a quantity that goes to zero when the quasirandomness parameter ${D}$ goes to infinity), and “almost all ${g \in G}$” denotes “for all ${g}$ in a subset of ${G}$ of density ${1-o(1)}$“. As a corollary, if ${A,B,C}$ have positive density in ${G}$ (by which we mean that ${\mu(A)}$ is bounded away from zero, uniformly in the quasirandomness parameter ${D}$, and similarly for ${B,C}$), then (if the quasirandomness parameter ${D}$ is sufficiently large) we can find elements ${g, x \in G}$ such that ${g \in A}$, ${x \in B}$, ${gx \in C}$. In fact we can find approximately ${\mu(A)\mu(B)\mu(C) |G|^2}$ such pairs ${(g,x)}$. To put it another way: if we choose ${g,x}$ uniformly and independently at random from ${G}$, then the events ${g \in A}$, ${x \in B}$, ${gx \in C}$ are approximately independent (thus the random variable ${(g,x,gx) \in G^3}$ resembles a uniformly distributed random variable on ${G^3}$ in some weak sense). One can also express this mixing property in integral form as

$\displaystyle \int_G \int_G f_1(g) f_2(x) f_3(gx)\ d\mu(g) d\mu(x) \approx (\int_G f_1\ d\mu) (\int_G f_2\ d\mu) (\int_G f_3\ d\mu)$

for any bounded functions ${f_1,f_2,f_3: G \rightarrow {\bf R}}$. (Of course, with ${G}$ being finite, one could replace the integrals here by finite averages if desired.) Or in probabilistic language, we have

$\displaystyle \mathop{\bf E} f_1(g) f_2(x) f_3(gx) \approx \mathop{\bf E} f_1(x_1) f_2(x_2) f_3(x_3)$

where ${g, x, x_1, x_2, x_3}$ are drawn uniformly and independently at random from ${G}$.

As observed in Gowers’ paper, one can iterate this observation to find “parallelopipeds” of any given dimension in dense subsets of ${G}$. For instance, applying (1) with ${A,B,C}$ replaced by ${A \cap hB}$, ${C \cap hD}$, and ${E \cap hF}$ one can assert (after some relabeling) that for ${g,h,x}$ chosen uniformly and independently at random from ${G}$, the events ${g \in A}$, ${h \in B}$, ${gh \in C}$, ${x \in D}$, ${gx \in E}$, ${hx \in F}$, ${ghx \in H}$ are approximately independent whenever ${A,B,C,D,E,F,H}$ are dense subsets of ${G}$; thus the tuple ${(g,h,gh,x,gh,hx,ghx)}$ resebles a uniformly distributed random variable in ${G^7}$ in some weak sense.

However, there are other tuples for which the above iteration argument does not seem to apply. One of the simplest tuples in this vein is the tuple ${(g, x, xg, gx)}$ in ${G^4}$, when ${g, x}$ are drawn uniformly at random from a quasirandom group ${G}$. Here, one does not expect the tuple to behave as if it were uniformly distributed in ${G^4}$, because there is an obvious constraint connecting the last two components ${gx, xg}$ of this tuple: they must lie in the same conjugacy class! In particular, if ${A}$ is a subset of ${G}$ that is the union of conjugacy classes, then the events ${gx \in A}$, ${xg \in A}$ are perfectly correlated, so that ${\mu( gx \in A, xg \in A)}$ is equal to ${\mu(A)}$ rather than ${\mu(A)^2}$. Our main result, though, is that in a quasirandom group, this is (approximately) the only constraint on the tuple. More precisely, we have

Theorem 1 Let ${G}$ be a ${D}$-quasirandom group, and let ${g, x}$ be drawn uniformly at random from ${G}$. Then for any ${f_1,f_2,f_3,f_4: G \rightarrow [-1,1]}$, we have

$\displaystyle \mathop{\bf E} f_1(g) f_2(x) f_3(gx) f_4(xg) = \mathop{\bf E} f_1(x_1) f_2(x_2) f_3(x_3) f_4(x_4) + o(1)$

where ${o(1)}$ goes to zero as ${D \rightarrow \infty}$, ${x_1,x_2,x_3}$ are drawn uniformly and independently at random from ${G}$, and ${x_4}$ is drawn uniformly at random from the conjugates of ${x_3}$ for each fixed choice of ${x_1,x_2,x_3}$.

This is the probabilistic formulation of the above theorem; one can also phrase the theorem in other formulations (such as an integral formulation), and this is detailed in the paper. This theorem leads to a number of recurrence results; for instance, as a corollary of this result, we have

$\displaystyle \mu(A) \mu(B)^2 - o(1) \leq \mu( A \cap gB \cap Bg ) \leq \mu(A) \mu(B) + o(1)$

for almost all ${g \in G}$, and any dense subsets ${A, B}$ of ${G}$; the lower and upper bounds are sharp, with the lower bound being attained when ${B}$ is randomly distributed, and the upper bound when ${B}$ is conjugation-invariant.

To me, the more interesting thing here is not the result itself, but how it is proven. Vitaly and I were not able to find a purely finitary way to establish this mixing theorem. Instead, we had to first use the machinery of ultraproducts (as discussed in this previous post) to convert the finitary statement about a quasirandom group to an infinitary statement about a type of infinite group which we call an ultra quasirandom group (basically, an ultraproduct of increasingly quasirandom finite groups). This is analogous to how the Furstenberg correspondence principle is used to convert a finitary combinatorial problem into an infinitary ergodic theory problem.

Ultra quasirandom groups come equipped with a finite, countably additive measure known as Loeb measure ${\mu_G}$, which is very analogous to the Haar measure of a compact group, except that in the case of ultra quasirandom groups one does not quite have a topological structure that would give compactness. Instead, one has a slightly weaker structure known as a ${\sigma}$-topology, which is like a topology except that open sets are only closed under countable unions rather than arbitrary ones. There are some interesting measure-theoretic and topological issues regarding the distinction between topologies and ${\sigma}$-topologies (and between Haar measure and Loeb measure), but for this post it is perhaps best to gloss over these issues and pretend that ultra quasirandom groups ${G}$ come with a Haar measure. One can then recast Theorem 1 as a mixing theorem for the left and right actions of the ultra approximate group ${G}$ on itself, which roughly speaking is the assertion that

$\displaystyle \int_G f_1(x) L_g f_2(x) L_g R_g f_3(x)\ d\mu_G(x) \approx 0 \ \ \ \ \ (2)$

for “almost all” ${g \in G}$, if ${f_1, f_2, f_3}$ are bounded measurable functions on ${G}$, with ${f_3}$ having zero mean on all conjugacy classes of ${G}$, where ${L_g, R_g}$ are the left and right translation operators

$\displaystyle L_g f(x) := f(g^{-1} x); \quad R_g f(x) := f(xg).$

To establish this mixing theorem, we use the machinery of idempotent ultrafilters, which is a particularly useful tool for understanding the ergodic theory of actions of countable groups ${G}$ that need not be amenable; in the non-amenable setting the classical ergodic averages do not make much sense, but ultrafilter-based averages are still available. To oversimplify substantially, the idempotent ultrafilter arguments let one establish mixing estimates of the form (2) for “many” elements ${g}$ of an infinite-dimensional parallelopiped known as an IP system (provided that the actions ${L_g,R_g}$ of this IP system obey some technical mixing hypotheses, but let’s ignore that for sake of this discussion). The claim then follows by using the quasirandomness hypothesis to show that if the estimate (2) failed for a large set of ${g \in G}$, then this large set would then contain an IP system, contradicting the previous claim.

Idempotent ultrafilters are an extremely infinitary type of mathematical object (one has to use Zorn’s lemma no fewer than three times just to construct one of these objects!). So it is quite remarkable that they can be used to establish a finitary theorem such as Theorem 1, though as is often the case with such infinitary arguments, one gets absolutely no quantitative control whatsoever on the error terms ${o(1)}$ appearing in that theorem. (It is also mildly amusing to note that our arguments involve the use of ultrafilters in two completely different ways: firstly in order to set up the ultraproduct that converts the finitary mixing problem to an infinitary one, and secondly to solve the infinitary mixing problem. Despite some superficial similarities, there appear to be no substantial commonalities between these two usages of ultrafilters.) There is already a fair amount of literature on using idempotent ultrafilter methods in infinitary ergodic theory, and perhaps by further development of ultraproduct correspondence principles, one can use such methods to obtain further finitary consequences (although the state of the art for idempotent ultrafilter ergodic theory has not advanced much beyond the analysis of two commuting shifts ${L_g, R_g}$ currently, which is the main reason why our arguments only handle the pattern ${(g,x,xg,gx)}$ and not more sophisticated patterns).

We also have some miscellaneous other results in the paper. It turns out that by using the triangle removal lemma from graph theory, one can obtain a recurrence result that asserts that whenever ${A}$ is a dense subset of a finite group ${G}$ (not necessarily quasirandom), then there are ${\gg |G|^2}$ pairs ${(x,g)}$ such that ${x, gx, xg}$ all lie in ${A}$. Using a hypergraph generalisation of the triangle removal lemma known as the hypergraph removal lemma, one can obtain more complicated versions of this statement; for instance, if ${A}$ is a dense subset of ${G^2}$, then one can find ${\gg |G|^2}$ triples ${(x,y,g)}$ such that ${(x,y), (gx, y), (gx, gy), (gxg^{-1}, gyg^{-1})}$ all lie in ${A}$. But the method is tailored to the specific types of patterns given here, and we do not have a general method for obtaining recurrence or mixing properties for arbitrary patterns of words in some finite alphabet such as ${g,x,y}$.

We also give some properties of a model example of an ultra quasirandom group, namely the ultraproduct ${SL_2(F)}$ of ${SL_2(F_{p_n})}$ where ${p_n}$ is a sequence of primes going off to infinity. Thanks to the substantial recent progress (by Helfgott, Bourgain, Gamburd, Breuillard, and others) on understanding the expansion properties of the finite groups ${SL_2(F_{p_n})}$, we have a fair amount of knowledge on the ultraproduct ${SL_2(F)}$ as well; for instance any two elements of ${SL_2(F)}$ will almost surely generate a spectral gap. We don’t have any direct application of this particular ultra quasirandom group, but it might be interesting to study it further.

Given a function ${f: X \rightarrow Y}$ between two sets ${X, Y}$, we can form the graph

$\displaystyle \Sigma := \{ (x,f(x)): x\in X \},$

which is a subset of the Cartesian product ${X \times Y}$.

There are a number of “closed graph theorems” in mathematics which relate the regularity properties of the function ${f}$ with the closure properties of the graph ${\Sigma}$, assuming some “completeness” properties of the domain ${X}$ and range ${Y}$. The most famous of these is the closed graph theorem from functional analysis, which I phrase as follows:

Theorem 1 (Closed graph theorem (functional analysis)) Let ${X, Y}$ be complete normed vector spaces over the reals (i.e. Banach spaces). Then a function ${f: X \rightarrow Y}$ is a continuous linear transformation if and only if the graph ${\Sigma := \{ (x,f(x)): x \in X \}}$ is both linearly closed (i.e. it is a linear subspace of ${X \times Y}$) and topologically closed (i.e. closed in the product topology of ${X \times Y}$).

I like to think of this theorem as linking together qualitative and quantitative notions of regularity preservation properties of an operator ${f}$; see this blog post for further discussion.

The theorem is equivalent to the assertion that any continuous linear bijection ${f: X \rightarrow Y}$ from one Banach space to another is necessarily an isomorphism in the sense that the inverse map is also continuous and linear. Indeed, to see that this claim implies the closed graph theorem, one applies it to the projection from ${\Sigma}$ to ${X}$, which is a continuous linear bijection; conversely, to deduce this claim from the closed graph theorem, observe that the graph of the inverse ${f^{-1}}$ is the reflection of the graph of ${f}$. As such, the closed graph theorem is a corollary of the open mapping theorem, which asserts that any continuous linear surjection from one Banach space to another is open. (Conversely, one can deduce the open mapping theorem from the closed graph theorem by quotienting out the kernel of the continuous surjection to get a bijection.)

It turns out that there is a closed graph theorem (or equivalent reformulations of that theorem, such as an assertion that bijective morphisms between sufficiently “complete” objects are necessarily isomorphisms, or as an open mapping theorem) in many other categories in mathematics as well. Here are some easy ones:

Theorem 2 (Closed graph theorem (linear algebra)) Let ${X, Y}$ be vector spaces over a field ${k}$. Then a function ${f: X \rightarrow Y}$ is a linear transformation if and only if the graph ${\Sigma := \{ (x,f(x)): x \in X \}}$ is linearly closed.

Theorem 3 (Closed graph theorem (group theory)) Let ${X, Y}$ be groups. Then a function ${f: X \rightarrow Y}$ is a group homomorphism if and only if the graph ${\Sigma := \{ (x,f(x)): x \in X \}}$ is closed under the group operations (i.e. it is a subgroup of ${X \times Y}$).

Theorem 4 (Closed graph theorem (order theory)) Let ${X, Y}$ be totally ordered sets. Then a function ${f: X \rightarrow Y}$ is monotone increasing if and only if the graph ${\Sigma := \{ (x,f(x)): x \in X \}}$ is totally ordered (using the product order on ${X \times Y}$).

Remark 1 Similar results to the above three theorems (with similarly easy proofs) hold for other algebraic structures, such as rings (using the usual product of rings), modules, algebras, or Lie algebras, groupoids, or even categories (a map between categories is a functor iff its graph is again a category). (ADDED IN VIEW OF COMMENTS: further examples include affine spaces and ${G}$-sets (sets with an action of a given group ${G}$).) There are also various approximate versions of this theorem that are useful in arithmetic combinatorics, that relate the property of a map ${f}$ being an “approximate homomorphism” in some sense with its graph being an “approximate group” in some sense. This is particularly useful for this subfield of mathematics because there are currently more theorems about approximate groups than about approximate homomorphisms, so that one can profitably use closed graph theorems to transfer results about the former to results about the latter.

A slightly more sophisticated result in the same vein:

Theorem 5 (Closed graph theorem (point set topology)) Let ${X, Y}$ be compact Hausdorff spaces. Then a function ${f: X \rightarrow Y}$ is continuous if and only if the graph ${\Sigma := \{ (x,f(x)): x \in X \}}$ is topologically closed.

Indeed, the “only if” direction is easy, while for the “if” direction, note that if ${\Sigma}$ is a closed subset of ${X \times Y}$, then it is compact Hausdorff, and the projection map from ${\Sigma}$ to ${X}$ is then a bijective continuous map between compact Hausdorff spaces, which is then closed, thus open, and hence a homeomorphism, giving the claim.

Note that the compactness hypothesis is necessary: for instance, the function ${f: {\bf R} \rightarrow {\bf R}}$ defined by ${f(x) := 1/x}$ for ${x \neq 0}$ and ${f(0)=0}$ for ${x=0}$ is a function which has a closed graph, but is discontinuous.

A similar result (but relying on a much deeper theorem) is available in algebraic geometry, as I learned after asking this MathOverflow question:

Theorem 6 (Closed graph theorem (algebraic geometry)) Let ${X, Y}$ be normal projective varieties over an algebraically closed field ${k}$ of characteristic zero. Then a function ${f: X \rightarrow Y}$ is a regular map if and only if the graph ${\Sigma := \{ (x,f(x)): x \in X \}}$ is Zariski-closed.

Proof: (Sketch) For the only if direction, note that the map ${x \mapsto (x,f(x))}$ is a regular map from the projective variety ${X}$ to the projective variety ${X \times Y}$ and is thus a projective morphism, hence is proper. In particular, the image ${\Sigma}$ of ${X}$ under this map is Zariski-closed.

Conversely, if ${\Sigma}$ is Zariski-closed, then it is also a projective variety, and the projection ${(x,y) \mapsto x}$ is a projective morphism from ${\Sigma}$ to ${X}$, which is clearly quasi-finite; by the characteristic zero hypothesis, it is also separated. Applying (Grothendieck’s form of) Zariski’s main theorem, this projection is the composition of an open immersion and a finite map. As projective varieties are complete, the open immersion is an isomorphism, and so the projection from ${\Sigma}$ to ${X}$ is finite. Being injective and separable, the degree of this finite map must be one, and hence ${k(\Sigma)}$ and ${k(X)}$ are isomorphic, hence (by normality of ${X}$) ${k[\Sigma]}$ is contained in (the image of) ${k[X]}$, which makes the map from ${X}$ to ${\Sigma}$ regular, which makes ${f}$ regular. $\Box$

The counterexample of the map ${f: k \rightarrow k}$ given by ${f(x) := 1/x}$ for ${x \neq 0}$ and ${f(0) := 0}$ demonstrates why the projective hypothesis is necessary. The necessity of the normality condition (or more precisely, a weak normality condition) is demonstrated by (the projective version of) the map ${(t^2,t^3) \mapsto t}$ from the cusipdal curve ${\{ (t^2,t^3): t \in k \}}$ to ${k}$. (If one restricts attention to smooth varieties, though, normality becomes automatic.) The necessity of characteristic zero is demonstrated by (the projective version of) the inverse of the Frobenius map ${x \mapsto x^p}$ on a field ${k}$ of characteristic ${p}$.

There are also a number of closed graph theorems for topological groups, of which the following is typical (see Exercise 3 of these previous blog notes):

Theorem 7 (Closed graph theorem (topological group theory)) Let ${X, Y}$ be ${\sigma}$-compact, locally compact Hausdorff groups. Then a function ${X \rightarrow Y}$ is a continuous homomorphism if and only if the graph ${\Sigma := \{ (x,f(x)): x \in X \}}$ is both group-theoretically closed and topologically closed.

The hypotheses of being ${\sigma}$-compact, locally compact, and Hausdorff can be relaxed somewhat, but I doubt that they can be eliminated entirely (though I do not have a ready counterexample for this).

In several complex variables, it is a classical theorem (see e.g. Lemma 4 of this blog post) that a holomorphic function from a domain in ${{\bf C}^n}$ to ${{\bf C}^n}$ is locally injective if and only if it is a local diffeomorphism (i.e. its derivative is everywhere non-singular). This leads to a closed graph theorem for complex manifolds:

Theorem 8 (Closed graph theorem (complex manifolds)) Let ${X, Y}$ be complex manifolds. Then a function ${f: X \rightarrow Y}$ is holomorphic if and only if the graph ${\Sigma := \{ (x,f(x)): x \in X \}}$ is a complex manifold (using the complex structure inherited from ${X \times Y}$) of the same dimension as ${X}$.

Indeed, one applies the previous observation to the projection from ${\Sigma}$ to ${X}$. The dimension requirement is needed, as can be seen from the example of the map ${f: {\bf C} \rightarrow {\bf C}}$ defined by ${f(z) =1/z}$ for ${z \neq 0}$ and ${f(0)=0}$.

(ADDED LATER:) There is a real analogue to the above theorem:

Theorem 9 (Closed graph theorem (real manifolds)) Let ${X, Y}$ be real manifolds. Then a function ${f: X \rightarrow Y}$ is continuous if and only if the graph ${\Sigma := \{ (x,f(x)): x \in X \}}$ is a real manifold of the same dimension as ${X}$.

This theorem can be proven by applying invariance of domain (discussed in this previous post) to the projection of ${\Sigma}$ to ${X}$, to show that it is open if ${\Sigma}$ has the same dimension as ${X}$.

Note though that the analogous claim for smooth real manifolds fails: the function ${f: {\bf R} \rightarrow {\bf R}}$ defined by ${f(x) := x^{1/3}}$ has a smooth graph, but is not itself smooth.

(ADDED YET LATER:) Here is an easy closed graph theorem in the symplectic category:

Theorem 10 (Closed graph theorem (symplectic geometry)) Let ${X = (X,\omega_X)}$ and ${Y = (Y,\omega_Y)}$ be smooth symplectic manifolds of the same dimension. Then a smooth map ${f: X \rightarrow Y}$ is a symplectic morphism (i.e. ${f^* \omega_Y = \omega_X}$) if and only if the graph ${\Sigma := \{(x,f(x)): x \in X \}}$ is a Lagrangian submanifold of ${X \times Y}$ with the symplectic form ${\omega_X \oplus -\omega_Y}$.

In view of the symplectic rigidity phenomenon, it is likely that the smoothness hypotheses on ${f,X,Y}$ can be relaxed substantially, but I will not try to formulate such a result here.

There are presumably many further examples of closed graph theorems (or closely related theorems, such as criteria for inverting a morphism, or open mapping type theorems) throughout mathematics; I would be interested to know of further examples.

$\Box$

This is a sequel to my previous blog post “Cayley graphs and the geometry of groups“. In that post, the concept of a Cayley graph of a group ${G}$ was used to place some geometry on that group ${G}$. In this post, we explore a variant of that theme, in which (fragments of) a Cayley graph on ${G}$ is used to describe the basic algebraic structure of ${G}$, and in particular on elementary word identities in ${G}$. Readers who are familiar with either category theory or group homology/cohomology will recognise these concepts lurking not far beneath the surface; we wil remark briefly on these connections later in this post. However, no knowledge of categories or cohomology is needed for the main discussion, which is primarily focused on elementary group theory.

Throughout this post, we fix a single group ${G = (G,\cdot)}$, which is allowed to be non-abelian and/or infinite. All our graphs will be directed, with loops and multiple edges permitted.

In the previous post, we drew the entire Cayley graph of a group ${G}$. Here, we will be working much more locally, and will only draw the portions of the Cayley graph that are relevant to the discussion. In this graph, the vertices are elements ${x}$ of the group ${G}$, and one draws a directed edge from ${x}$ to ${xg}$ labeled (or “coloured”) by the group element ${g}$ for any ${x, g \in G}$; the graph consisting of all such vertices and edges will be denoted ${Cay(G,G)}$. Thus, a typical edge in ${Cay(G,G)}$ looks like this:

Figure 1.

One usually does not work with the complete Cayley graph ${Cay(G,G)}$. It is customary to instead work with smaller Cayley graphs ${Cay(G,S)}$, in which the edge colours ${g}$ are restricted to a smaller subset of ${G}$, such as a set of generators for ${G}$. As we will be working locally, we will in fact work with even smaller fragments of ${Cay(G,G)}$ at a time; in particular, we only use a handful of colours (no more than nine, in fact, for any given diagram), and we will not require these colours to generate the entire group (we do not care if the Cayley graph is connected or not, as this is a global property rather than a local one).

Cayley graphs are left-invariant: for any ${a \in G}$, the left translation map ${x \mapsto ax}$ is a graph isomorphism. To emphasise this left invariance, we will usually omit the vertex labels, and leave only the coloured directed edge, like so:

Figure 2.

This is analogous to how, in undergraduate mathematics and physics, vectors in Euclidean space are often depicted as arrows of a given magnitude and direction, with the initial and final points of this arrow being of secondary importance only. (Indeed, this depiction of vectors in a vector space can be viewed as an abelian special case of the more general depiction of group elements used in this post.)

Let us define a diagram to be a finite directed graph ${H = (V,E)}$, with edges coloured by elements of ${G}$, which has at least one graph homomorphism into the complete Cayley graph ${Cay(G,G)}$ of ${G}$; thus there exists a map ${\phi: V \rightarrow G}$ (not necessarily injective) with the property that ${\phi(w) = \phi(v) g}$ whenever ${(v,w)}$ is a directed edge in ${H}$ coloured by a group element ${g \in G}$. Informally, a diagram is a finite subgraph of a Cayley graph with the vertex labels omitted, and with distinct vertices permitted to represent the same group element. Thus, for instance, the single directed edge displayed in Figure 2 is a very simple example of a diagram. An even simpler example of a diagram would be a depiction of the identity element:

Figure 3.

We will however omit the identity loops in our diagrams in order to reduce clutter.

We make the obvious remark that any directed edge in a diagram can be coloured by at most one group element ${g}$, since ${y=xg, y=xh}$ implies ${g=h}$. This simple observation provides a way to prove group theoretic identities using diagrams: to show that two group elements ${g, h}$ are equal, it suffices to show that they connect together (with the same orientation) the same pair of vertices in a diagram.

Remark 1 One can also interpret these diagrams as commutative diagrams in a category in which all the objects are copies of ${G}$, and the morphisms are right-translation maps. However, we will deviate somewhat from the category theoretic way of thinking here by focusing on the geometric arrangement and shape of these diagrams, rather than on their abstract combinatorial description. In particular, we view the arrows more as distorted analogues of vector arrows, than as the abstract arrows appearing in category theory.

Just as vector addition can be expressed via concatenation of arrows, group multiplication can be described by concatenation of directed edges. Indeed, for any ${x,g,h \in G}$, the vertices ${x, xg, xgh}$ can be connected by the following triangular diagram:

Figure 4.

In a similar spirit, inversion is described by the following diagram:

Figure 5.

We make the pedantic remark though that we do not consider a ${g^{-1}}$ edge to be the reversal of the ${g}$ edge, but rather as a distinct edge that just happens to have the same initial and final endpoints as the reversal of the ${g}$ edge. (This will be of minor importance later, when we start integrating “${1}$-forms” on such edges.)

A fundamental operation for us will be that of gluing two diagrams together.

Lemma 1 ((Labeled) gluing) Let ${D_1 = (V_1,E_1), D_2 = (V_2,E_2)}$ be two diagrams of a given group ${G}$. Suppose that the intersection ${D_1 \cap D_2 := (V_1 \cap V_2, E_1 \cap E_2)}$ of the two diagrams connects all of ${V_1 \cap V_2}$ (i.e. any two elements of ${V_1 \cap V_2}$ are joined by a path in ${D_1 \cap D_2}$). Then the union ${D_1 \cup D_2 := (V_1 \cup V_2, E_1 \cup E_2)}$ is also a diagram of ${G}$.

Proof: By hypothesis, we have graph homomorphisms ${\phi_1: D_1 \rightarrow Cay(G,G)}$, ${\phi_2: D_2 \rightarrow Cay(G,G)}$. If they agree on ${D_1 \cap D_2}$ then one simply glues together the two homomorphisms to create a new graph homomorphism ${\phi: D_1 \cup D_2 \rightarrow Cay(G,G)}$. If they do not agree, one can apply a left translation to either ${\phi_1}$ or ${\phi_2}$ to make the two diagrams agree on at least one vertex of ${D_1 \cap D_2}$; then by the connected nature of ${D_1 \cap D_2}$ we see that they now must agree on all vertices of ${D_1 \cap D_2}$, and then we can form the glued graph homomorphism as before. $\Box$

The above lemma required one to specify the label the vertices of ${D_1,D_2}$ (in order to form the intersection ${D_1 \cap D_2}$ and union ${D_1 \cup D_2}$). However, if one is presented with two diagrams ${D_1, D_2}$ with unlabeled vertices, one can identify some partial set of vertices of ${D_1}$ with a partial set of vertices of ${D_2}$ of matching cardinality. Provided that the subdiagram common to ${D_1}$ and ${D_2}$ after this identification connects all of the common vertices together, we may use the above lemma to create a glued diagram ${D}$.

For instance, if a diagram ${D}$ contains two of the three edges in the triangular diagram in Figure 4, one can “fill in” the triangle by gluing in the third edge:

Figure 6.

One can use glued diagrams to demonstrate various basic group-theoretic identities. For instance, by gluing together two copies of the triangular diagram in Figure 4 to create the glued diagram

Figure 7.

and then filling in two more triangles, we obtain a tetrahedral diagram that demonstrates the associative law ${(gh)k = g(hk)}$:

Figure 8.

Similarly, by gluing together two copies of Figure 4 with three copies of Figure 5 in an appropriate order, we can demonstrate the Abel identity ${(gh)^{-1} = h^{-1} g^{-1}}$:

Figure 9.

In addition to gluing, we will also use the trivial operation of erasing: if ${D}$ is a diagram for a group ${G}$, then any subgraph of ${D}$ (formed by removing vertices and/or edges) is also a diagram of ${G}$. This operation is not strictly necessary for our applications, but serves to reduce clutter in the pictures.

If two group elements ${g, h}$ commute, then we obtain a parallelogram as a diagram, exactly as in the vector space case:

Figure 10.

In general, of course, two arbitrary group elements ${g,h}$ will fail to commute, and so this parallelogram is no longer available. However, various substitutes for this diagram exist. For instance, if we introduce the conjugate ${g^h := h^{-1} g h}$ of one group element ${g}$ by another, then we have the following slightly distorted parallelogram:

Figure 11.

By appropriate gluing and filling, this can be used to demonstrate the homomorphism properties of a conjugation map ${g \mapsto g^h}$:

Figure 12.

Figure 13.

Another way to replace the parallelogram in Figure 10 is to introduce the commutator ${[g,h] := g^{-1}h^{-1}gh}$ of two elements, in which case we can perturb the parallelogram into a pentagon:

Figure 14.

We will tend to depict commutator edges as being somewhat shorter than the edges generating that commutator, reflecting a “perturbative” or “nilpotent” philosophy. (Of course, to fully reflect a nilpotent perspective, one should orient commutator edges in a different dimension from their generating edges, but of course the diagrams drawn here do not have enough dimensions to display this perspective easily.) We will also be adopting a “Lie” perspective of interpreting groups as behaving like perturbations of vector spaces, in particular by trying to draw all edges of the same colour as being approximately (though not perfectly) parallel to each other (and with approximately the same length).

Gluing the above pentagon with the conjugation parallelogram and erasing some edges, we discover a “commutator-conjugate” triangle, describing the basic identity ${g^h = g [g,h]}$:

Figure 15.

Other gluings can also give the basic relations between commutators and conjugates. For instance, by gluing the pentagon in Figure 14 with its reflection, we see that ${[g,h] = [h,g]^{-1}}$. The following diagram, obtained by gluing together copies of Figures 11 and 15, demonstrates that ${[h,g^{-1}] = [g,h]^{g^{-1}}}$,

Figure 16.

while this figure demonstrates that ${[g,hk] = [g,k] [g,h]^k}$:

Figure 17.

Now we turn to a more sophisticated identity, the Hall-Witt identity

$\displaystyle [[g,h],k^g] [[k,g],h^k] [[h,k],g^h] = 1,$

which is the fully noncommutative version of the more well-known Jacobi identity for Lie algebras.

The full diagram for the Hall-Witt identity resembles a slightly truncated parallelopiped. Drawing this truncated paralleopiped in full would result in a rather complicated looking diagram, so I will instead display three components of this diagram separately, and leave it to the reader to mentally glue these three components back to form the full parallelopiped. The first component of the diagram is formed by gluing together three pentagons from Figure 14, and looks like this:

Figure 18.

This should be thought of as the “back” of the truncated parallelopiped needed to establish the Hall-Witt identity.

While it is not needed for proving the Hall-Witt identity, we also observe for future reference that we may also glue in some distorted parallelograms and obtain a slightly more complicated diagram:

Figure 19.

To form the second component, let us now erase all interior components of Figure 18 or Figure 19:

Figure 20.

Then we fill in three distorted parallelograms:

Figure 21.

This is the second component, and is the “front” of the truncated praallelopiped, minus the portions exposed by the truncation.

Finally, we turn to the third component. We begin by erasing the outer edges from the second component in Figure 21:

Figure 22.

We glue in three copies of the commutator-conjugate triangle from Figure 15:

Figure 23.

But now we observe that we can fill in three pentagons, and obtain a small triangle with edges ${[[g,h],k^g] [[k,g],h^k] [[h,k],g^h]}$:

Figure 24.

Erasing everything except this triangle gives the Hall-Witt identity. Alternatively, one can glue together Figures 18, 21, and 24 to obtain a truncated parallelopiped which one can view as a geometric representation of the proof of the Hall-Witt identity.

Among other things, I found these diagrams to be useful to visualise group cohomology; I give a simple example of this below, developing an analogue of the Hall-Witt identity for ${2}$-cocycles.

This is an addendum to last quarter’s course notes on Hilbert’s fifth problem, which I am in the process of reviewing in order to transcribe them into a book (as was done similarly for several other sets of lecture notes on this blog). When reviewing the zeroth set of notes in particular, I found that I had made a claim (Proposition 11 from those notes) which asserted, roughly speaking, that any sufficiently large nilprogression was an approximate group, and promised to prove it later in the course when we had developed the ability to calculate efficiently in nilpotent groups. As it turned out, I managed finish the course without the need to develop these calculations, and so the proposition remained unproven. In order to rectify this, I will use this post to lay out some of the basic algebra of nilpotent groups, and use it to prove the above proposition, which turns out to be a bit tricky. (In my paper with Breuillard and Green, we avoid the need for this proposition by restricting attention to a special type of nilprogression, which we call a nilprogression in ${C}$-normal form, for which the computations are simpler.)

There are several ways to think about nilpotent groups; for instance one can use the model example of the Heisenberg group

$\displaystyle H(R) :=\begin{pmatrix} 1 & R & R \\ 0 & 1 & R\\ 0 & 0 & 1 \end{pmatrix}$

over an arbitrary ring ${R}$ (which need not be commutative), or more generally any matrix group consisting of unipotent upper triangular matrices, and view a general nilpotent group as being an abstract generalisation of such concrete groups. (In the case of nilpotent Lie groups, at least, this is quite an accurate intuition, thanks to Engel’s theorem.) Or, one can adopt a Lie-theoretic viewpoint and try to think of nilpotent groups as somehow arising from nilpotent Lie algebras; this intuition is rigorous when working with nilpotent Lie groups (at least when the characteristic is large, in order to avoid issues coming from the denominators in the Baker-Campbell-Hausdorff formula), but also retains some conceptual value in the non-Lie setting. In particular, nilpotent groups (particularly finitely generated ones) can be viewed in some sense as “nilpotent Lie groups over ${{\bf Z}}$“, even though Lie theory does not quite work perfectly when the underlying scalars merely form an integral domain instead of a field.

Another point of view, which arises naturally both in analysis and in algebraic geometry, is to view nilpotent groups as modeling “infinitesimal” perturbations of the identity, where the infinitesimals have a certain finite order. For instance, given a (not necessarily commutative) ring ${R}$ without identity (representing all the “small” elements of some larger ring or algebra), we can form the powers ${R^j}$ for ${j=1,2,\ldots}$, defined as the ring generated by ${j}$-fold products ${r_1 \ldots r_j}$ of elements ${r_1,\ldots,r_j}$ in ${R}$; this is an ideal of ${R}$ which represents the elements which are “${j^{th}}$ order” in some sense. If one then formally adjoins an identity ${1}$ onto the ring ${R}$, then for any ${s \geq 1}$, the multiplicative group ${G := 1+R \hbox{ mod } R^{s+1}}$ is a nilpotent group of step at most ${s}$. For instance, if ${R}$ is the ring of strictly upper ${s \times s}$ matrices (over some base ring), then ${R^{s+1}}$ vanishes and ${G}$ becomes the group of unipotent upper triangular matrices over the same ring, thus recovering the previous matrix-based example. In analysis applications, ${R}$ might be a ring of operators which are somehow of “order” ${O(\epsilon)}$ or ${O(\hbar)}$ for some small parameter ${\epsilon}$ or ${\hbar}$, and one wishes to perform Taylor expansions up to order ${O(\epsilon^s)}$ or ${O(\hbar^s)}$, thus discarding (i.e. quotienting out) all errors in ${R^{s+1}}$.

From a dynamical or group-theoretic perspective, one can also view nilpotent groups as towers of central extensions of a trivial group. Finitely generated nilpotent groups can also be profitably viewed as a special type of polycylic group; this is the perspective taken in this previous blog post. Last, but not least, one can view nilpotent groups from a combinatorial group theory perspective, as being words from some set of generators of various “degrees” subject to some commutation relations, with commutators of two low-degree generators being expressed in terms of higher degree objects, and all commutators of a sufficiently high degree vanishing. In particular, generators of a given degree can be moved freely around a word, as long as one is willing to generate commutator errors of higher degree.

With this last perspective, in particular, one can start computing in nilpotent groups by adopting the philosophy that the lowest order terms should be attended to first, without much initial concern for the higher order errors generated in the process of organising the lower order terms. Only after the lower order terms are in place should attention then turn to higher order terms, working successively up the hierarchy of degrees until all terms are dealt with. This turns out to be a relatively straightforward philosophy to implement in many cases (particularly if one is not interested in explicit expressions and constants, being content instead with qualitative expansions of controlled complexity), but the arguments are necessarily recursive in nature and as such can become a bit messy, and require a fair amount of notation to express precisely. So, unfortunately, the arguments here will be somewhat cumbersome and notation-heavy, even if the underlying methods of proof are relatively simple.

In this final set of course notes, we discuss how (a generalisation of) the expansion results obtained in the preceding notes can be used for some nnumber-theoretic applications, and in particular to locate almost primes inside orbits of thin groups, following the work of Bourgain, Gamburd, and Sarnak. We will not attempt here to obtain the sharpest or most general results in this direction, but instead focus on the simplest instances of these results which are still illustrative of the ideas involved.

One of the basic general problems in analytic number theory is to locate tuples of primes of a certain form; for instance, the famous (and still unsolved) twin prime conjecture asserts that there are infinitely many pairs ${(n_1,n_2)}$ in the line ${\{ (n_1,n_2) \in {\bf Z}^2: n_2-n_1=2\}}$ in which both entries are prime. In a similar spirit, one of the Landau conjectures (also still unsolved) asserts that there are infinitely many primes in the set ${\{ n^2+1: n \in {\bf Z} \}}$. The Mersenne prime conjecture (also unsolved) asserts that there are infinitely many primes in the set ${\{ 2^n - 1: n \in {\bf Z} \}}$, and so forth.

More generally, given some explicit subset ${V}$ in ${{\bf R}^d}$ (or ${{\bf C}^d}$, if one wishes), such as an algebraic variety, one can ask the question of whether there are infinitely many integer lattice points ${(n_1,\ldots,n_d)}$ in ${V \cap {\bf Z}^d}$ in which all the coefficients ${n_1,\ldots,n_d}$ are simultaneously prime; let us refer to such points as prime points.

At this level of generality, this problem is impossibly difficult. Indeed, even the much simpler problem of deciding whether the set ${V \cap {\bf Z}^d}$ is non-empty (let alone containing prime points) when ${V}$ is a hypersurface ${\{ x \in {\bf R}^d: P(x) = 0 \}}$ cut out by a polynomial ${P}$ is essentially Hilbert’s tenth problem, which is known to be undecidable in general by Matiyasevich’s theorem. So one needs to restrict attention to a more special class of sets ${V}$, in which the question of finding integer points is not so difficult. One model case is to consider orbits ${V = \Gamma b}$, where ${b \in {\bf Z}^d}$ is a fixed lattice vector and ${\Gamma}$ is some discrete group that acts on ${{\bf Z}^d}$ somehow (e.g. ${\Gamma}$ might be embedded as a subgroup of the special linear group ${SL_d({\bf Z})}$, or on the affine group ${SL_d({\bf Z}) \ltimes {\bf Z}^d}$). In such a situation it is then quite easy to show that ${V = V \cap {\bf Z}^d}$ is large; for instance, ${V}$ will be infinite precisely when the stabiliser of ${b}$ in ${\Gamma}$ has infinite index in ${\Gamma}$.

Even in this simpler setting, the question of determining whether an orbit ${V = \Gamma b}$ contains infinitely prime points is still extremely difficult; indeed the three examples given above of the twin prime conjecture, Landau conjecture, and Mersenne prime conjecture are essentially of this form (possibly after some slight modification of the underlying ring ${{\bf Z}}$, see this paper of Bourgain-Gamburd-Sarnak for details), and are all unsolved (and generally considered well out of reach of current technology). Indeed, the list of non-trivial orbits ${V = \Gamma b}$ which are known to contain infinitely many prime points is quite slim; Euclid’s theorem on the infinitude of primes handles the case ${V = {\bf Z}}$, Dirichlet’s theorem handles infinite arithmetic progressions ${V = a{\bf Z} + r}$, and a somewhat complicated result of Green, Tao, and Ziegler handles “non-degenerate” affine lattices in ${{\bf Z}^d}$ of rank two or more (such as the lattice of length ${d}$ arithmetic progressions), but there are few other positive results known that are not based on the above cases (though we will note the remarkable theorem of Friedlander and Iwaniec that there are infinitely many primes of the form ${a^2+b^4}$, and the related result of Heath-Brown that there are infinitely many primes of the form ${a^3+2b^3}$, as being in a kindred spirit to the above results, though they are not explicitly associated to an orbit of a reasonable action as far as I know).

On the other hand, much more is known if one is willing to replace the primes by the larger set of almost primes – integers with a small number of prime factors (counting multiplicity). Specifically, for any ${r \geq 1}$, let us call an ${r}$-almost prime an integer which is the product of at most ${r}$ primes, and possibly by the unit ${-1}$ as well. Many of the above sorts of questions which are open for primes, are known for ${r}$-almost primes for ${r}$ sufficiently large. For instance, with regards to the twin prime conjecture, it is a result of Chen that there are infinitely many pairs ${p,p+2}$ where ${p}$ is a prime and ${p+2}$ is a ${2}$-almost prime; in a similar vein, it is a result of Iwaniec that there are infinitely many ${2}$-almost primes of the form ${n^2+1}$. On the other hand, it is still open for any fixed ${r}$ whether there are infinitely many Mersenne numbers ${2^n-1}$ which are ${r}$-almost primes. (For the superficially similar situation with the numbers ${2^n+1}$, it is in fact believed (but again unproven) that there are only finitely many ${r}$-almost primes for any fixed ${r}$ (the Fermat prime conjecture).)

The main tool that allows one to count almost primes in orbits is sieve theory. The reason for this lies in the simple observation that in order to ensure that an integer ${n}$ of magnitude at most ${x}$ is an ${r}$-almost prime, it suffices to guarantee that ${n}$ is not divisible by any prime less than ${x^{1/(r+1)}}$. Thus, to create ${r}$-almost primes, one can start with the integers up to some large threshold ${x}$ and remove (or “sieve out”) all the integers that are multiples of any prime ${p}$ less than ${x^{1/(r+1)}}$. The difficulty is then to ensure that a sufficiently non-trivial quantity of integers remain after this process, for the purposes of finding points in the given set ${V}$.

The most basic sieve of this form is the sieve of Eratosthenes, which when combined with the inclusion-exclusion principle gives the Legendre sieve (or exact sieve), which gives an exact formula for quantities such as the number ${\pi(x,z)}$ of natural numbers less than or equal to ${x}$ that are not divisible by any prime less than or equal to a given threshold ${z}$. Unfortunately, when one tries to evaluate this formula, one encounters error terms which grow exponentially in ${z}$, rendering this sieve useful only for very small thresholds ${z}$ (of logarithmic size in ${x}$). To improve the sieve level up to a small power of ${x}$ such as ${x^{1/(r+1)}}$, one has to replace the exact sieve by upper bound sieves and lower bound sieves which only seek to obtain upper or lower bounds on quantities such as ${\pi(x,z)}$, but contain a polynomial number of terms rather than an exponential number. There are a variety of such sieves, with the two most common such sieves being combinatorial sieves (such as the beta sieve), based on various combinatorial truncations of the inclusion-exclusion formula, and the Selberg upper bound sieve, based on upper bounds that are the square of a divisor sum. (There is also the large sieve, which is somewhat different in nature and based on ${L^2}$ almost orthogonality considerations, rather than on any actual sieving, to obtain upper bounds.) We will primarily work with a specific sieve in this notes, namely the beta sieve, and we will not attempt to optimise all the parameters of this sieve (which ultimately means that the almost primality parameter ${r}$ in our results will be somewhat large). For a more detailed study of sieve theory, see the classic text of Halberstam and Richert, or the more recent texts of Iwaniec-Kowalski or of Friedlander-Iwaniec.

Very roughly speaking, the end result of sieve theory is that excepting some degenerate and “exponentially thin” settings (such as those associated with the Mersenne primes), all the orbits which are expected to have a large number of primes, can be proven to at least have a large number of ${r}$-almost primes for some finite ${r}$. (Unfortunately, there is a major obstruction, known as the parity problem, which prevents sieve theory from lowering ${r}$ all the way to ${1}$; see this blog post for more discussion.) One formulation of this principle was established by Bourgain, Gamburd, and Sarnak:

Theorem 1 (Bourgain-Gamburd-Sarnak) Let ${\Gamma}$ be a subgroup of ${SL_2({\bf Z})}$ which is not virtually solvable. Let ${f: {\bf Z}^4 \rightarrow {\bf Z}}$ be a polynomial with integer coefficients obeying the following primitivity condition: for any positive integer ${q}$, there exists ${A \in \Gamma \subset {\bf Z}^4}$ such that ${f(A)}$ is coprime to ${q}$. Then there exists an ${r \geq 1}$ such that there are infinitely many ${A \in \Gamma}$ with ${f(A)}$ non-zero and ${r}$-almost prime.

This is not the strongest version of the Bourgain-Gamburd-Sarnak theorem, but it captures the general flavour of their results. Note that the theorem immediately implies an analogous result for orbits ${\Gamma b \subset {\bf Z}^2}$, in which ${f}$ is now a polynomial from ${{\bf Z}^2}$ to ${{\bf Z}}$, and one uses ${f(Ab)}$ instead of ${f(A)}$. It is in fact conjectured that one can set ${r=1}$ here, but this is well beyond current technology. For the purpose of reaching ${r=1}$, it is very natural to impose the primitivity condition, but as long as one is content with larger values of ${r}$, it is possible to relax the primitivity condition somewhat; see the paper of Bourgain, Gamburd, and Sarnak for more discussion.

By specialising to the polynomial ${f: \begin{pmatrix} a & b \\ c & d \end{pmatrix} \rightarrow abcd}$, we conclude as a corollary that as long as ${\Gamma}$ is primitive in the sense that it contains matrices with all coefficients coprime to ${q}$ for any given ${q}$, then ${\Gamma}$ contains infinitely many matrices whose elements are all ${r}$-almost primes for some ${r}$ depending only on ${\Gamma}$. For further applications of these sorts of results, for instance to Appolonian packings, see the paper of Bourgain, Gamburd, and Sarnak.

It turns out that to prove Theorem 1, the Cayley expansion results in ${SL_2(F_p)}$ from the previous set of notes are not quite enough; one needs a more general Cayley expansion result in ${SL_2({\bf Z}/q{\bf Z})}$ where ${q}$ is square-free but not necessarily prime. The proof of this expansion result uses the same basic methods as in the ${SL_2(F_p)}$ case, but is significantly more complicated technically, and we will only discuss it briefly here. As such, we do not give a complete proof of Theorem 1, but hopefully the portion of the argument presented here is still sufficient to give an impression of the ideas involved.

In the last three notes, we discussed the Bourgain-Gamburd expansion machine and two of its three ingredients, namely quasirandomness and product theorems, leaving only the non-concentration ingredient to discuss. We can summarise the results of the last three notes, in the case of fields of prime order, as the following theorem.

Theorem 1 (Non-concentration implies expansion in ${SL_d}$) Let ${p}$ be a prime, let ${d \geq 1}$, and let ${S}$ be a symmetric set of elements in ${G := SL_d(F_p)}$ of cardinality ${|S|=k}$ not containing the identity. Write ${\mu := \frac{1}{|S|} \sum_{s\in S}\delta_s}$, and suppose that one has the non-concentration property

$\displaystyle \sup_{H < G}\mu^{(n)}(H) < |G|^{-\kappa} \ \ \ \ \ (1)$

for some ${\kappa>0}$ and some even integer ${n \leq \Lambda \log |G|}$. Then ${Cay(G,S)}$ is a two-sided ${\epsilon}$-expander for some ${\epsilon>0}$ depending only on ${k, d, \kappa,\Lambda}$.

Proof: From (1) we see that ${\mu^{(n)}}$ is not supported in any proper subgroup ${H}$ of ${G}$, which implies that ${S}$ generates ${G}$. The claim now follows from the Bourgain-Gamburd expansion machine (Theorem 2 of Notes 4), the product theorem (Theorem 1 of Notes 5), and quasirandomness (Exercise 8 of Notes 3). $\Box$

Remark 1 The same argument also works if we replace ${F_p}$ by the field ${F_{p^j}}$ of order ${p^j}$ for some bounded ${j}$. However, there is a difficulty in the regime when ${j}$ is unbounded, because the quasirandomness property becomes too weak for the Bourgain-Gamburd expansion machine to be directly applicable. On theother hand, the above type of theorem was generalised to the setting of cyclic groups ${{\bf Z}/q{\bf Z}}$ with ${q}$ square-free by Varju, to arbitrary ${q}$ by Bourgain and Varju, and to more general algebraic groups than ${SL_d}$ and square-free ${q}$ by Salehi Golsefidy and Varju. It may be that some modification of the proof techniques in these papers may also be able to handle the field case ${F_{p^j}}$ with unbounded ${j}$.

It thus remains to construct tools that can establish the non-concentration property (1). The situation is particularly simple in ${SL_2(F_p)}$, as we have a good understanding of the subgroups of that group. Indeed, from Theorem 14 from Notes 5, we obtain the following corollary to Theorem 1:

Corollary 2 (Non-concentration implies expansion in ${SL_2}$) Let ${p}$ be a prime, and let ${S}$ be a symmetric set of elements in ${G := SL_2(F_p)}$ of cardinality ${|S|=k}$ not containing the identity. Write ${\mu := \frac{1}{|S|} \sum_{s\in S}\delta_s}$, and suppose that one has the non-concentration property

$\displaystyle \sup_{B}\mu^{(n)}(B) < |G|^{-\kappa} \ \ \ \ \ (2)$

for some ${\kappa>0}$ and some even integer ${n \leq \Lambda \log |G|}$, where ${B}$ ranges over all Borel subgroups of ${SL_2(\overline{F})}$. Then, if ${|G|}$ is sufficiently large depending on ${k,\kappa,\Lambda}$, ${Cay(G,S)}$ is a two-sided ${\epsilon}$-expander for some ${\epsilon>0}$ depending only on ${k, \kappa,\Lambda}$.

It turns out (2) can be verified in many cases by exploiting the solvable nature of the Borel subgroups ${B}$. We give two examples of this in these notes. The first result, due to Bourgain and Gamburd (with earlier partial results by Gamburd and by Shalom) generalises Selberg’s expander construction to the case when ${S}$ generates a thin subgroup of ${SL_2({\bf Z})}$:

Theorem 3 (Expansion in thin subgroups) Let ${S}$ be a symmetric subset of ${SL_2({\bf Z})}$ not containing the identity, and suppose that the group ${\langle S \rangle}$ generated by ${S}$ is not virtually solvable. Then as ${p}$ ranges over all sufficiently large primes, the Cayley graphs ${Cay(SL_2(F_p), \pi_p(S))}$ form a two-sided expander family, where ${\pi_p: SL_2({\bf Z}) \rightarrow SL_2(F_p)}$ is the usual projection.

Remark 2 One corollary of Theorem 3 (or of the non-concentration estimate (3) below) is that ${\pi_p(S)}$ generates ${SL_2(F_p)}$ for all sufficiently large ${p}$, if ${\langle S \rangle}$ is not virtually solvable. This is a special case of a much more general result, known as the strong approximation theorem, although this is certainly not the most direct way to prove such a theorem. Conversely, the strong approximation property is used in generalisations of this result to higher rank groups than ${SL_2}$.

Exercise 1 In the converse direction, if ${\langle S\rangle}$ is virtually solvable, show that for sufficiently large ${p}$, ${\pi_p(S)}$ fails to generate ${SL_2(F_p)}$. (Hint: use Theorem 14 from Notes 5 to prevent ${SL_2(F_p)}$ from having bounded index solvable subgroups.)

Exercise 2 (Lubotzsky’s 1-2-3 problem) Let ${S := \{ \begin{pmatrix}1 & \pm 3 \\ 0 & 1 \end{pmatrix}, \begin{pmatrix}1 & 0 \\ \pm 3 & 1 \end{pmatrix}}$.

• (i) Show that ${S}$ generates a free subgroup of ${SL_2({\bf Z})}$. (Hint: use a ping-pong argument, as in Exercise 23 of Notes 2.)
• (ii) Show that if ${v, w}$ are two distinct elements of the sector ${\{ (x,y) \in {\bf R}^2_+: x/2 < y < 2x \}}$, then there os no element ${g \in \langle S \rangle}$ for which ${gv = w}$. (Hint: this is another ping-pong argument.) Conclude that ${\langle S \rangle}$ has infinite index in ${SL_2({\bf Z})}$. (Contrast this with the situation in which the ${3}$ coefficients in ${S}$ are replaced by ${1}$ or ${2}$, in which case ${\langle S \rangle}$ is either all of ${SL_2({\bf Z})}$, or a finite index subgroup, as demonstrated in Exercise 23 of Notes 2).
• (iii) Show that ${Cay(SL_2(F_p), \pi_p(S))}$ for sufficiently large primes ${p}$ form a two-sided expander family.

Remark 3 Theorem 3 has been generalised to arbitrary linear groups, and with ${F_p}$ replaced by ${{\bf Z}/q{\bf Z}}$ for square-free ${q}$; see this paper of Salehi Golsefidy and Varju. In this more general setting, the condition of virtual solvability must be replaced by the condition that the connected component of the Zariski closure of ${\langle S \rangle}$ is perfect. An effective version of Theorem 3 (with completely explicit constants) was recently obtained by Kowalski.

The second example concerns Cayley graphs constructed using random elements of ${SL_2(F_p)}$.

Theorem 4 (Random generators expand) Let ${p}$ be a prime, and let ${x,y}$ be two elements of ${SL_2(F_p)}$ chosen uniformly at random. Then with probability ${1-o_{p \rightarrow \infty}(1)}$, ${Cay(SL_2(F_p), \{x,x^{-1},y,y^{-1}\})}$ is a two-sided ${\epsilon}$-expander for some absolute constant ${\epsilon}$.

Remark 4 As with Theorem 3, Theorem 4 has also been extended to a number of other groups, such as the Suzuki groups (in this paper of Breuillard, Green, and Tao), and more generally to finite simple groups of Lie type of bounded rank (in forthcoming work of Breuillard, Green, Guralnick, and Tao). There are a number of other constructions of expanding Cayley graphs in such groups (and in other interesting groups, such as the alternating groups) beyond those discussed in these notes; see this recent survey of Lubotzky for further discussion. It has been conjectured by Lubotzky and Weiss that any pair ${x,y}$ of (say) ${SL_2(F_p)}$ that generates the group, is a two-sided ${\epsilon}$-expander for an absolute constant ${\epsilon}$: in the case of ${SL_2(F_p)}$, this has been established for a density one set of primes by Breuillard and Gamburd.

— 1. Expansion in thin subgroups —

We now prove Theorem 3. The first observation is that the expansion property is monotone in the group ${\langle S \rangle}$:

Exercise 3 Let ${S, S'}$ be symmetric subsets of ${SL_2({\bf Z})}$ not containing the identity, such that ${\langle S \rangle \subset \langle S' \rangle}$. Suppose that ${Cay(SL_2(F_p), \pi_p(S))}$ is a two-sided expander family for sufficiently large primes ${p}$. Show that ${Cay(SL_2(F_p), \pi_p(S'))}$ is also a two-sided expander family.

As a consequence, Theorem 3 follows from the following two statments:

Theorem 5 (Tits alternative) Let ${\Gamma \subset SL_2({\bf Z})}$ be a group. Then exactly one of the following statements holds:

• (i) ${\Gamma}$ is virtually solvable.
• (ii) ${\Gamma}$ contains a copy of the free group ${F_2}$ of two generators as a subgroup.

Theorem 6 (Expansion in free groups) Let ${x,y \in SL_2({\bf Z})}$ be generators of a free subgroup of ${SL_2({\bf Z})}$. Then as ${p}$ ranges over all sufficiently large primes, the Cayley graphs ${Cay(SL_2(F_p), \pi_p(\{x,y,x^{-1},y^{-1}\}))}$ form a two-sided expander family.

Theorem 5 is a special case of the famous Tits alternative, which among other things allows one to replace ${SL_2({\bf Z})}$ by ${GL_d(k)}$ for any ${d \geq 1}$ and any field ${k}$ of characteristic zero (and fields of positive characteristic are also allowed, if one adds the requirement that ${\Gamma}$ be finitely generated). We will not prove the full Tits alternative here, but instead just give an ad hoc proof of the special case in Theorem 5 in the following exercise.

Exercise 4 Given any matrix ${g \in SL_2({\bf Z})}$, the singular values are ${\|g\|_{op}}$ and ${\|g\|_{op}^{-1}}$, and we can apply the singular value decomposition to decompose

$\displaystyle g = u_1(g) \|g\|_{op} v_1^*(g) + u_2(g) \|g\|_{op}^{-1} v_2(g)^*$

where ${u_1(g),u_2(g)\in {\bf C}^2}$ and ${v_1(g), v_2(g) \in {\bf C}^2}$ are orthonormal bases. (When ${\|g\|_{op}>1}$, these bases are uniquely determined up to phase rotation.) We let ${\tilde u_1(g) \in {\bf CP}^1}$ be the projection of ${u_1(g)}$ to the projective complex plane, and similarly define ${\tilde v_2(g)}$.

Let ${\Gamma}$ be a subgroup of ${SL_2({\bf Z})}$. Call a pair ${(u,v) \in {\bf CP}^1 \times {\bf CP}^1}$ a limit point of ${\Gamma}$ if there exists a sequence ${g_n \in \Gamma}$ with ${\|g_n\|_{op} \rightarrow \infty}$ and ${(\tilde u_1(g_n), \tilde v_2(g_n)) \rightarrow (u,v)}$.

• (i) Show that if ${\Gamma}$ is infinite, then there is at least one limit point.
• (ii) Show that if ${(u,v)}$ is a limit point, then so is ${(v,u)}$.
• (iii) Show that if there are two limit points ${(u,v), (u',v')}$ with ${\{u,v\} \cap \{u',v'\} = \emptyset}$, then there exist ${g,h \in \Gamma}$ that generate a free group. (Hint: Choose ${(\tilde u_1(g), \tilde v_2(g))}$ close to ${(u,v)}$ and ${(\tilde u_1(h),\tilde v_2(h))}$ close to ${(u',v')}$, and consider the action of ${g}$ and ${h}$ on ${{\bf CP}^1}$, and specifically on small neighbourhoods of ${u,v,u',v'}$, and set up a ping-pong type situation.)
• (iv) Show that if ${g \in SL_2({\bf Z})}$ is hyperbolic (i.e. it has an eigenvalue greater than 1), with eigenvectors ${u,v}$, then the projectivisations ${(\tilde u,\tilde v)}$ of ${u,v}$ form a limit point. Similarly, if ${g}$ is regular parabolic (i.e. it has an eigenvalue at 1, but is not the identity) with eigenvector ${u}$, show that ${(\tilde u,\tilde bu)}$ is a limit point.
• (v) Show that if ${\Gamma}$ has no free subgroup of two generators, then all hyperbolic and regular parabolic elements of ${\Gamma}$ have a common eigenvector. Conclude that all such elements lie in a solvable subgroup of ${\Gamma}$.
• (vi) Show that if an element ${g \in SL_2({\bf Z})}$ is neither hyperbolic nor regular parabolic, and is not a multiple of the identity, then ${g}$ is conjugate to a rotation by ${\pi/2}$ (in particular, ${g^2=-1}$).
• (vii) Establish Theorem 5. (Hint: show that two square roots of ${-1}$ in ${SL_2({\bf Z})}$ cannot multiply to another square root of ${-1}$.)

Now we prove Theorem 6. Let ${\Gamma}$ be a free subgroup of ${SL_2({\bf Z})}$ generated by two generators ${x,y}$. Let ${\mu := \frac{1}{4} (\delta_x +\delta_{x^{-1}} + \delta_y + \delta_{y^{-1}})}$ be the probability measure generating a random walk on ${SL_2({\bf Z})}$, thus ${(\pi_p)_* \mu}$ is the corresponding generator on ${SL_2(F_p)}$. By Corollary 2, it thus suffices to show that

$\displaystyle \sup_{B}((\pi_p)_* \mu)^{(n)}(B) < p^{-\kappa} \ \ \ \ \ (3)$

for all sufficiently large ${p}$, some absolute constant ${\kappa>0}$, and some even ${n = O(\log p)}$ (depending on ${p}$, of course), where ${B}$ ranges over Borel subgroups.

As ${\pi_p}$ is a homomorphism, one has ${((\pi_p)_* \mu)^{(n)}(B) = (\pi_p)_* (\mu^{(n)})(B) = \mu^{(n)}(\pi_p^{-1}(B))}$ and so it suffices to show that

$\displaystyle \sup_{B} \mu^{(n)}(\pi_p^{-1}(B)) < p^{-\kappa}.$

To deal with the supremum here, we will use an argument of Bourgain and Gamburd, taking advantage of the fact that all Borel groups of ${SL_2}$ obey a common group law, the point being that free groups such as ${\Gamma}$ obey such laws only very rarely. More precisely, we use the fact that the Borel groups are solvable of derived length two; in particular we have

$\displaystyle [[a,b],[c,d]] = 1 \ \ \ \ \ (4)$

for all ${a,b,c,d \in B}$. Now, ${\mu^{(n)}}$ is supported on matrices in ${SL_2({\bf Z})}$ whose coefficients have size ${O(\exp(O(n)))}$ (where we allow the implied constants to depend on the choice of generators ${x,y}$), and so ${(\pi_p)_*( \mu^{(n)} )}$ is supported on matrices in ${SL_2(F_p)}$ whose coefficients also have size ${O(\exp(O(n)))}$. If ${n}$ is less than a sufficiently small multiple of ${\log p}$, these coefficients are then less than ${p^{1/10}}$ (say). As such, if ${\tilde a,\tilde b,\tilde c,\tilde d \in SL_2({\bf Z})}$ lie in the support of ${\mu^{(n)}}$ and their projections ${a = \pi_p(\tilde a), \ldots, d = \pi_p(\tilde d)}$ obey the word law (4) in ${SL_2(F_p)}$, then the original matrices ${\tilde a, \tilde b, \tilde c, \tilde d}$ obey the word law (4) in ${SL_2({\bf Z})}$. (This lifting of identities from the characteristic ${p}$ setting of ${SL_2(F_p)}$ to the characteristic ${0}$ setting of ${SL_2({\bf Z})}$ is a simple example of the “Lefschetz principle”.)

To summarise, if we let ${E_{n,p,B}}$ be the set of all elements of ${\pi_p^{-1}(B)}$ that lie in the support of ${\mu^{(n)}}$, then (4) holds for all ${a,b,c,d \in E_{n,p,B}}$. This severely limits the size of ${E_{n,p,B}}$ to only be of polynomial size, rather than exponential size:

Proposition 7 Let ${E}$ be a subset of the support of ${\mu^{(n)}}$ (thus, ${E}$ consists of words in ${x,y,x^{-1},y^{-1}}$ of length ${n}$) such that the law (4) holds for all ${a,b,c,d \in E}$. Then ${|E| \ll n^2}$.

The proof of this proposition is laid out in the exercise below.

Exercise 5 Let ${\Gamma}$ be a free group generated by two generators ${x,y}$. Let ${B}$ be the set of all words of length at most ${n}$ in ${x,y,x^{-1},y^{-1}}$.

• (i) Show that if ${a,b \in \Gamma}$ commute, then ${a, b}$ lie in the same cyclic group, thus ${a = c^i, b = c^j}$ for some ${c \in \Gamma}$ and ${i,j \in {\bf Z}}$.
• (ii) Show that if ${a \in \Gamma}$, there are at most ${O(n)}$ elements of ${B}$ that commute with ${a}$.
• (iii) Show that if ${a,c \in \Gamma}$, there are at most ${O(n)}$ elements ${b}$ of ${B}$ with ${[a,b] = c}$.
• (iv) Prove Proposition 7.

Now we can conclude the proof of Theorem 3:

Exercise 6 Let ${\Gamma}$ be a free group generated by two generators ${x,y}$.

• (i) Show that ${\| \mu^{(n)} \|_{\ell^\infty(\Gamma)} \ll c^n}$ for some absolute constant ${0 < c<1}$. (For much more precise information on ${\mu^{(n)}}$, see this paper of Kesten.)
• (ii) Conclude the proof of Theorem 3.

— 2. Random generators expand —

We now prove Theorem 4. Let ${{\bf F}_2}$ be the free group on two formal generators ${a,b}$, and let ${\mu := \frac{1}{4}(\delta_a + \delta_b + \delta_{a^{-1}}+ \delta_{b^{-1}}}$ be the generator of the random walk. For any word ${w \in {\bf F}_2}$ and any ${x,y}$ in a group ${G}$, let ${w(x,y) \in G}$ be the element of ${G}$ formed by substituting ${x,y}$ for ${a,b}$ respectively in the word ${w}$; thus ${w}$ can be viewed as a map ${w: G \times G \rightarrow G}$ for any group ${G}$. Observe that if ${w}$ is drawn randomly using the distribution ${\mu^{(n)}}$, and ${x,y \in SL_2(F_p)}$, then ${w(x,y)}$ is distributed according to the law ${\tilde \mu^{(n)}}$, where ${\tilde \mu := \frac{1}{4}(\delta_x + \delta_y + \delta_{x^{-1}}+ \delta_{y^{-1}})}$. Applying Corollary 2, it suffices to show that whenever ${p}$ is a large prime and ${x,y}$ are chosen uniformly and independently at random from ${SL_2(F_p)}$, that with probability ${1-o_{p \rightarrow \infty}(1)}$, one has

$\displaystyle \sup_B {\bf P}_w ( w(x,y) \in B ) \leq p^{-\kappa} \ \ \ \ \ (5)$

for some absolute constant ${\kappa}$, where ${B}$ ranges over all Borel subgroups of ${SL_2(\overline{F_p})}$ and ${w}$ is drawn from the law ${\mu^{(n)}}$ for some even natural number ${n = O(\log p)}$.

Let ${B_n}$ denote the words in ${{\bf F}_2}$ of length at most ${n}$. We may use the law (4) to obtain good bound on the supremum in (5) assuming a certain non-degeneracy property of the word evaluations ${w(x,y)}$:

Exercise 7 Let ${n}$ be a natural number, and suppose that ${x,y \in SL_2(F_p)}$ is such that ${w(x,y) \neq 1}$ for ${w \in B_{100n} \backslash \{1\}}$. Show that

$\displaystyle \sup_B {\bf P}_w ( w(x,y) \in B ) \ll \exp(-cn)$

for some absolute constant ${c>0}$, where ${w}$ is drawn from the law ${\mu^{(n)}}$. (Hint: use (4) and the hypothesis to lift the problem up to ${{\bf F}_2}$, at which point one can use Proposition 7 and Exercise 6.)

In view of this exercise, it suffices to show that with probability ${1-o_{p \rightarrow\infty}(1)}$, one has ${w(x,y) \neq 1}$ for all ${w \in B_{100n} \backslash \{1\}}$ for some ${n}$ comparable to a small multiple of ${\log p}$. As ${B_{100n}}$ has ${\exp(O(n))}$ elements, it thus suffices by the union bound to show that

$\displaystyle {\bf P}_{x,y}(w(x,y)=1) \leq p^{-\gamma} \ \ \ \ \ (6)$

for some absolute constant ${\gamma > 0}$, and any ${w \in {\bf F}_2 \backslash \{1\}}$ of length less than ${c\log p}$ for some sufficiently small absolute constant ${c>0}$.

Let us now fix a non-identity word ${w}$ of length ${|w|}$ less than ${c\log p}$, and consider ${w}$ as a function from ${SL_2(k) \times SL_2(k)}$ to ${SL_2(k)}$ for an arbitrary field ${k}$. We can identify ${SL_2(k)}$ with the set ${\{ (a,b,c,d)\in k^4: ad-bc=1\}}$. A routine induction then shows that the expression ${w((a,b,c,d),(a',b',c',d'))}$ is then a polynomial in the eight variables ${a,b,c,d,a',b',c',d'}$ of degree ${O(|w|)}$ and coefficients which are integers of size ${O( \exp( O(|w|) ) )}$. Let us then make the additional restriction to the case ${a,a' \neq 0}$, in which case we can write ${d = \frac{bc+1}{a}}$ and ${d' =\frac{b'c'+1}{a'}}$. Then ${w((a,b,c,d),(a',b',c',d'))}$ is now a rational function of ${a,b,c,a',b',c'}$ whose numerator is a polynomial of degree ${O(|w|)}$ and coefficients of size ${O( \exp( O(|w|) ) )}$, and the denominator is a monomial of ${a,a'}$ of degree ${O(|w|)}$.

We then specialise this rational function to the field ${k=F_p}$. It is conceivable that when one does so, the rational function collapses to the constant polynomial ${(1,0,0,1)}$, thus ${w((a,b,c,d),(a',b',c',d'))=1}$ for all ${(a,b,c,d),(a',b',c',d') \in SL_2(F_p)}$ with ${a,a' \neq 0}$. (For instance, this would be the case if ${w(x,y) = x^{|SL_2(F_p)|}}$, by Lagrange’s theorem, if it were not for the fact that ${|w|}$ is far too large here.) But suppose that this rational function does not collapse to the constant rational function. Applying the Schwarz-Zippel lemma (Exercise 23 from Notes 5), we then see that the set of pairs ${(a,b,c,d),(a',b',c',d') \in SL_2(F_p)}$ with ${a,a' \neq 0}$ and ${w((a,b,c,d),(a',b',c',d'))=1}$ is at most ${O( |w| p^5 )}$; adding in the ${a=0}$ and ${a'=0}$ cases, one still obtains a bound of ${O(|w|p^5)}$, which is acceptable since ${|SL_2(F_p)|^2 \sim p^6}$ and ${|w| = O( \log p )}$. Thus, the only remaining case to consider is when the rational function ${w((a,b,c,d),(a',b',c',d'))}$ is identically ${1}$ on ${SL_2(F_p)}$ with ${a,a' \neq 0}$.

Now we perform another “Lefschetz principle” maneuvre to change the underlying field. Recall that the denominator of rational function ${w((a,b,c,d),(a',b',c',d'))}$ is monomial in ${a,a'}$, and the numerator has coefficients of size ${O(\exp(O(|w|)))}$. If ${|w|}$ is less than ${c\log p}$ for a sufficiently small ${p}$, we conclude in particular (for ${p}$ large enough) that the coefficients all have magnitude less than ${p}$. As such, the only way that this function can be identically ${1}$ on ${SL_2(F_p)}$ is if it is identically ${1}$ on ${SL_2(k)}$ for all ${k}$ with ${a,a' \neq 0}$, and hence for ${a=0}$ or ${a'=0}$ also by taking Zariski closures.

On the other hand, we know that for some choices of ${k}$, e.g. ${k={\bf R}}$, ${SL_2(k)}$ contains a copy ${\Gamma}$ of the free group on two generators (see e.g. Exercise 23 of Notes 2). As such, it is not possible for any non-identity word ${w}$ to be identically trivial on ${SL_2(k) \times SL_2(k)}$. Thus this case cannot actually occur, completing the proof of (6) and hence of Theorem 4.

Remark 5 We see from the above argument that the existence of subgroups ${\Gamma}$ of an algebraic group with good “independence” properties – such as that of generating a free group – can be useful in studying the expansion properties of that algebraic group, even if the field of interest in the latter is distinct from that of the former. For more complicated algebraic groups than ${SL_2}$, in which laws such as (4) are not always available, it turns out to be useful to place further properties on the subgroup ${\Gamma}$, for instance by requiring that all non-abelian subgroups of that group be Zariski dense (a property which has been called strong density), as this turns out to be useful for preventing random walks from concentrating in proper algebraic subgroups. See this paper of Breuillard, Guralnick, Green and Tao for constructions of strongly dense free subgroups of algebraic groups and further discussion.

In the previous set of notes, we saw that one could derive expansion of Cayley graphs from three ingredients: non-concentration, product theorems, and quasirandomness. Quasirandomness was discussed in Notes 3. In the current set of notes, we discuss product theorems. Roughly speaking, these theorems assert that in certain circumstances, a finite subset ${A}$ of a group ${G}$ either exhibits expansion (in the sense that ${A^3}$, say, is significantly larger than ${A}$), or is somehow “close to” or “trapped” by a genuine group.

Theorem 1 (Product theorem in ${SL_d(k)}$) Let ${d \geq 2}$, let ${k}$ be a finite field, and let ${A}$ be a finite subset of ${G := SL_d(k)}$. Let ${\epsilon >0}$ be sufficiently small depending on ${d}$. Then at least one of the following statements holds:

• (Expansion) One has ${|A^3| \geq |A|^{1+\epsilon}}$.
• (Close to ${G}$) One has ${|A| \geq |G|^{1-O_d(\epsilon)}}$.
• (Trapping) ${A}$ is contained in a proper subgroup of ${G}$.

We will prove this theorem (which was proven first in the ${d=2,3}$ cases for fields ${F}$ of prime order by Helfgott, and then for ${d=2}$ and general ${F}$ by Dinai, and finally to general ${d}$ and ${F}$ independently by Pyber-Szabo and by Breuillard-Green-Tao) later in this notes. A more qualitative version of this proposition was also previously obtained by Hrushovski. There are also generalisations of the product theorem of importance to number theory, in which the field ${k}$ is replaced by a cyclic ring ${{\bf Z}/q{\bf Z}}$ (with ${q}$ not necessarily prime); this was achieved first for ${d=2}$ and ${q}$ square-free by Bourgain, Gamburd, and Sarnak, by Varju for general ${d}$ and ${q}$ square-free, and finally by this paper of Bourgain and Varju for arbitrary ${d}$ and ${q}$.

Exercise 1 (Girth bound) Assuming Theorem 1, show that whenever ${S}$ is a symmetric set of generators of ${SL_d(k)}$ for some finite field ${k}$ and some ${d\geq 2}$, then any element of ${SL_d(k)}$ can be expressed as the product of ${O_d( \log^{O_d(1)} |k| )}$ elements from ${S}$. (Equivalently, if we add the identity element to ${S}$, then ${S^m = SL_d(k)}$ for some ${m = O_d( \log^{O_d(1)} |k| )}$.) This is a special case of a conjecture of Babai and Seress, who conjectured that the bound should hold uniformly for all finite simple groups (in particular, the implied constants here should not actually depend on ${d}$. The methods used to handle the ${SL_d}$ case can handle other finite groups of Lie type of bounded rank, but at present we do not have bounds that are independent of the rank. On the other hand, a recent paper of Helfgott and Seress has almost resolved the conjecture for the permutation groups ${A_n}$.

A key tool to establish product theorems is an argument which is sometimes referred to as the pivot argument. To illustrate this argument, let us first discuss a much simpler (and older) theorem, essentially due to Freiman, which has a much weaker conclusion but is valid in any group ${G}$:

Theorem 2 (Baby product theorem) Let ${G}$ be a group, and let ${A}$ be a finite non-empty subset of ${G}$. Then one of the following statements hold:

• (Expansion) One has ${|A^{-1} A| \geq \frac{3}{2} |A|}$.
• (Close to a subgroup) ${A}$ is contained in a left-coset of a group ${H}$ with ${|H| < \frac{3}{2} |A|}$.

To prove this theorem, we suppose that the first conclusion does not hold, thus ${|A^{-1} A| <\frac{3}{2} |A|}$. Our task is then to place ${A}$ inside the left-coset of a fairly small group ${H}$.

To do this, we take a group element ${g \in G}$, and consider the intersection ${A\cap gA}$. A priori, the size of this set could range from anywhere from ${0}$ to ${|A|}$. However, we can use the hypothesis ${|A^{-1} A| < \frac{3}{2} |A|}$ to obtain an important dichotomy, reminiscent of the classical fact that two cosets ${gH, hH}$ of a subgroup ${H}$ of ${G}$ are either identical or disjoint:

Proposition 3 (Dichotomy) If ${g \in G}$, then exactly one of the following occurs:

• (Non-involved case) ${A \cap gA}$ is empty.
• (Involved case) ${|A \cap gA| > \frac{|A|}{2}}$.

Proof: Suppose we are not in the pivot case, so that ${A \cap gA}$ is non-empty. Let ${a}$ be an element of ${A \cap gA}$, then ${a}$ and ${g^{-1} a}$ both lie in ${A}$. The sets ${A^{-1} a}$ and ${A^{-1} g^{-1} a}$ then both lie in ${A^{-1} A}$. As these sets have cardinality ${|A|}$ and lie in ${A^{-1}A}$, which has cardinality less than ${\frac{3}{2}|A|}$, we conclude from the inclusion-exclusion formula that

$\displaystyle |A^{-1} a \cap A^{-1} g^{-1} a| > \frac{|A|}{2}.$

But the left-hand side is equal to ${|A \cap gA|}$, and the claim follows. $\Box$

The above proposition provides a clear separation between two types of elements ${g \in G}$: the “non-involved” elements, which have nothing to do with ${A}$ (in the sense that ${A \cap gA = \emptyset}$, and the “involved” elements, which have a lot to do with ${A}$ (in the sense that ${|A \cap gA| > |A|/2}$. The key point is that there is a significant “gap” between the non-involved and involved elements; there are no elements that are only “slightly involved”, in that ${A}$ and ${gA}$ intersect a little but not a lot. It is this gap that will allow us to upgrade approximate structure to exact structure. Namely,

Proposition 4 The set ${H}$ of involved elements is a finite group, and is equal to ${A A^{-1}}$.

Proof: It is clear that the identity element ${1}$ is involved, and that if ${g}$ is involved then so is ${g^{-1}}$ (since ${A \cap g^{-1} A = g^{-1}(A \cap gA)}$. Now suppose that ${g, h}$ are both involved. Then ${A \cap gA}$ and ${A\cap hA}$ have cardinality greater than ${|A|/2}$ and are both subsets of ${A}$, and so have non-empty intersection. In particular, ${gA \cap hA}$ is non-empty, and so ${A \cap g^{-1} hA}$ is non-empty. By Proposition 3, this makes ${g^{-1} h}$ involved. It is then clear that ${H}$ is a group.

If ${g \in A A^{-1}}$, then ${A \cap gA}$ is non-empty, and so from Proposition 3 ${g}$ is involved. Conversely, if ${g}$ is involved, then ${g \in A A^{-1}}$. Thus we have ${H = A A^{-1}}$ as claimed. In particular, ${H}$ is finite. $\Box$

Now we can quickly wrap up the proof of Theorem 2. By construction, ${A \cap gA| > |A|/2}$ for all ${g \in H}$,which by double counting shows that ${|H| < 2|A|}$. As ${H = A A^{-1}}$, we see that ${A}$ is contained in a right coset ${Hg}$ of ${H}$; setting ${H' := g^{-1} H g}$, we conclude that ${A}$ is contained in a left coset ${gH'}$ of ${H'}$. ${H'}$ is a conjugate of ${H}$, and so ${|H'| < 2|A|}$. If ${h \in H'}$, then ${A}$ and ${Ah}$ both lie in ${H'}$ and have cardinality ${|A|}$, so must overlap; and so ${h \in A A^{-1}}$. Thus ${A A^{-1} = H'}$, and so ${|H'| < \frac{3}{2} |A|}$, and Theorem 2 follows.

Exercise 2 Show that the constant ${3/2}$ in Theorem 2 cannot be replaced by any larger constant.

Exercise 3 Let ${A \subset G}$ be a finite non-empty set such that ${|A^2| < 2|A|}$. Show that ${AA^{-1}=A^{-1} A}$. (Hint: If ${ab^{-1} \in A A^{-1}}$, show that ${ab^{-1} = c^{-1} d}$ for some ${c,d \in A}$.)

Exercise 4 Let ${A \subset G}$ be a finite non-empty set such that ${|A^2| < \frac{3}{2} |A|}$. Show that there is a finite group ${H}$ with ${|H| < \frac{3}{2} |A|}$ and a group element ${g \in G}$ such that ${A \subset Hg \cap gH}$ and ${H = A A^{-1}}$.

Below the fold, we give further examples of the pivot argument in other group-like situations, including Theorem 2 and also the “sum-product theorem” of Bourgain-Katz-Tao and Bourgain-Glibichuk-Konyagin.

Emmanuel Breuillard, Ben Green and I have just uploaded to the arXiv the short paper “A nilpotent Freiman dimension lemma“, submitted to the special volume of the European Journal of Combinatorics in honour of Yahya Ould Hamidoune.  This paper is a nonabelian (or more precisely, nilpotent) variant of the following additive combinatorics lemma of Freiman:

Freiman’s lemma.  Let A be a finite subset of a Euclidean space with $|A+A| \leq K|A|$.  Then A is contained in an affine subspace of dimension at most ${}\lfloor K-1 \rfloor$.

This can be viewed as a “cheap” version of the more well known theorem of Freiman that places sets of small doubling in a torsion-free abelian group inside a generalised arithmetic progression.  The advantage here is that the bound on the dimension is extremely explicit.

Our main result is

Theorem.  Let A be a finite subset of a simply-connected nilpotent Lie group G which is a K-approximate group (i.e. A is symmetric, contains the identity, and $A \cdot A$ can be covered by up to K left translates of A.  Then A can be covered by at most $K^{2+29K^9}$ left-translates of a closed connected Lie subgroup of dimension at most $K^9$.

We remark that our previous paper established a similar result, in which the dimension bound was improved to $6\log_2 K$, but at the cost of worsening the covering number to $O_K(1)$, and with a much more complicated proof  (91 pages instead of 8). Furthermore, the bound on $O_K(1)$ is ineffective, due to the use of ultraproducts in the argument (though it is likely that some extremely lousy explicit bound could eventually be squeezed out of the argument by finitising everything).  Note that the step of the ambient nilpotent group G does not influence the final bounds in the theorem, although we do of course need this step to be finite.  A simple quotienting argument allows one to deduce a corollary of the above theorem in which the ambient group is assumed to be residually torsion-free nilpotent instead of being a simply connected nilpotent Lie group, but we omit the statement of this corollary here.

To motivate the proof of this theorem, let us first show a simple case of an argument of Gleason, which is very much in the spirit of Freiman’s lemma:

Gleason Lemma (special case).  Let $A$ be a finite symmetric subset of a Euclidean space, and let $0 = H_0 \subset H_1 \subset \ldots \subset H_k$ be a sequence of subspaces in this space, such that the sets $2A \cap H_i$ are strictly increasing in i for $i=0,\ldots,k$.  Then $|5A| \geq k|A|$, where $5A = A+A+A+A+A$.

Proof.    By hypothesis, for each $i=1,\ldots,k$, the projection $B_i$ of $2A \cap H_i$ to $H_i / H_{i-1}$ is non-trivial, finite, and symmetric.   In particular, since the vector space $H_i/H_{i-1}$ is torsion-free, $B_i+B_i$ is strictly larger than $B_i$.  Equivalently, one can find $a_i$ in $(2A \cap H_i) + (2A \cap H_i)$ that does not lie in $(2A \cap H_i) + H_{i-1}$; in particular, $a_i \in 4A \cap H_i$ and $a_i+A$ is disjoint from $H_{i-1}+A$.  As a consequence, the $a_i+A$ are disjoint and lie in 5A, whence the claim. $\Box$

Note that by combining the contrapositive of this lemma with a greedy algorithm, one can show that any K-approximate group in a Euclidean space is contained in a subspace of dimension at most $K^4$, which is a weak version of Freiman’s lemma.

To extend the argument to the nilpotent setting we use the following idea.  Observe that any non-trivial genuine subgroup H of a nilpotent group G will contain at least one non-trivial central element; indeed, by intersecting H with the lower central series $G = G_1 \geq G_2 \geq G_3 \geq \ldots$ of G, and considering the last intersection $H \cap G_k$ which is non-trivial, one obtains the claim.  It turns out that one can adapt this argument to approximate groups, so that any sufficiently large K-approximate subgroup A of G will contain a non-trivial element that centralises a large fraction of A.  Passing to this large fraction and quotienting out the central element, we obtain a new approximate group.    If, after a bounded number of steps, this procedure gives an approximate group of bounded size, we are basically done.  If, however, the process continues, then by using some Lie group theory, one can find a long sequence $H_0 \leq H_1 \leq \ldots \leq H_k$ of connected Lie subgroups of G, such that the sets $A^2 \cap H_i$ are strictly increasing in i.   Using some Lie group theory and the hypotheses on G, one can deduce that the group $\langle A^2 \cap H_{i+1}\rangle$ generated by $A^2 \cap H_{i+1}$ is much larger than $\langle A^2 \cap H_i \rangle$, in the sense that the latter group has infinite index in the former.   It then turns out that the Gleason argument mentioned above can be adapted to this setting.