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There are multiple purposes to this blog post.

The first purpose is to announce the uploading of the paper “New equidistribution estimates of Zhang type, and bounded gaps between primes” by D.H.J. Polymath, which is the main output of the Polymath8a project on bounded gaps between primes, to the arXiv, and to describe the main results of this paper below the fold.

The second purpose is to roll over the previous thread on all remaining Polymath8a-related matters (e.g. updates on the submission status of the paper) to a fresh thread. (Discussion of the ongoing Polymath8b project is however being kept on a separate thread, to try to reduce confusion.)

The final purpose of this post is to coordinate the writing of a retrospective article on the Polymath8 experience, which has been solicited for the Newsletter of the European Mathematical Society. I suppose that this could encompass both the Polymath8a and Polymath8b projects, even though the second one is still ongoing (but I think we will soon be entering the endgame there). I think there would be two main purposes of such a retrospective article. The first one would be to tell a story about the process of conducting mathematical research, rather than just describe the outcome of such research; this is an important aspect of the subject which is given almost no attention in most mathematical writing, and it would be good to be able to capture some sense of this process while memories are still relatively fresh. The other would be to draw some tentative conclusions with regards to what the strengths and weaknesses of a Polymath project are, and how appropriate such a format would be for other mathematical problems than bounded gaps between primes. In my opinion, the bounded gaps problem had some fairly unique features that made it particularly amenable to a Polymath project, such as (a) a high level of interest amongst the mathematical community in the problem; (b) a very focused objective (“improve {H}!”), which naturally provided an obvious metric to measure progress; (c) the modular nature of the project, which allowed for people to focus on one aspect of the problem only, and still make contributions to the final goal; and (d) a very reasonable level of ambition (for instance, we did not attempt to prove the twin prime conjecture, which in my opinion would make a terrible Polymath project at our current level of mathematical technology). This is not an exhaustive list of helpful features of the problem; I would welcome other diagnoses of the project by other participants.

With these two objectives in mind, I propose a format for the retrospective article consisting of a brief introduction to the polymath concept in general and the polymath8 project in particular, followed by a collection of essentially independent contributions by different participants on their own experiences and thoughts. Finally we could have a conclusion section in which we make some general remarks on the polymath project (such as the remarks above). I’ve started a dropbox subfolder for this article (currently in a very skeletal outline form only), and will begin writing a section on my own experiences; other participants are of course encouraged to add their own sections (it is probably best to create separate files for these, and then input them into the main file retrospective.tex, to reduce edit conflicts. If there are participants who wish to contribute but do not currently have access to the Dropbox folder, please email me and I will try to have you added (or else you can supply your thoughts by email, or in the comments to this post; we may have a section for shorter miscellaneous comments from more casual participants, for people who don’t wish to write a lengthy essay on the subject).

As for deadlines, the EMS Newsletter would like a submitted article by mid-April in order to make the June issue, but in the worst case, it will just be held over until the issue after that.

Read the rest of this entry »

I’ve just uploaded to the arXiv the paper “Finite time blowup for an averaged three-dimensional Navier-Stokes equation“, submitted to J. Amer. Math. Soc.. The main purpose of this paper is to formalise the “supercriticality barrier” for the global regularity problem for the Navier-Stokes equation, which roughly speaking asserts that it is not possible to establish global regularity by any “abstract” approach which only uses upper bound function space estimates on the nonlinear part of the equation, combined with the energy identity. This is done by constructing a modification of the Navier-Stokes equations with a nonlinearity that obeys essentially all of the function space estimates that the true Navier-Stokes nonlinearity does, and which also obeys the energy identity, but for which one can construct solutions that blow up in finite time. Results of this type had been previously established by Montgomery-Smith, Gallagher-Paicu, and Li-Sinai for variants of the Navier-Stokes equation without the energy identity, and by Katz-Pavlovic and by Cheskidov for dyadic analogues of the Navier-Stokes equations in five and higher dimensions that obeyed the energy identity (see also the work of Plechac and Sverak and of Hou and Lei that also suggest blowup for other Navier-Stokes type models obeying the energy identity in five and higher dimensions), but to my knowledge this is the first blowup result for a Navier-Stokes type equation in three dimensions that also obeys the energy identity. Intriguingly, the method of proof in fact hints at a possible route to establishing blowup for the true Navier-Stokes equations, which I am now increasingly inclined to believe is the case (albeit for a very small set of initial data).

To state the results more precisely, recall that the Navier-Stokes equations can be written in the form

\displaystyle  \partial_t u + (u \cdot \nabla) u = \nu \Delta u + \nabla p

for a divergence-free velocity field {u} and a pressure field {p}, where {\nu>0} is the viscosity, which we will normalise to be one. We will work in the non-periodic setting, so the spatial domain is {{\bf R}^3}, and for sake of exposition I will not discuss matters of regularity or decay of the solution (but we will always be working with strong notions of solution here rather than weak ones). Applying the Leray projection {P} to divergence-free vector fields to this equation, we can eliminate the pressure, and obtain an evolution equation

\displaystyle  \partial_t u = \Delta u + B(u,u) \ \ \ \ \ (1)

purely for the velocity field, where {B} is a certain bilinear operator on divergence-free vector fields (specifically, {B(u,v) = -\frac{1}{2} P( (u \cdot \nabla) v + (v \cdot \nabla) u)}. The global regularity problem for Navier-Stokes is then equivalent to the global regularity problem for the evolution equation (1).

An important feature of the bilinear operator {B} appearing in (1) is the cancellation law

\displaystyle  \langle B(u,u), u \rangle = 0

(using the {L^2} inner product on divergence-free vector fields), which leads in particular to the fundamental energy identity

\displaystyle  \frac{1}{2} \int_{{\bf R}^3} |u(T,x)|^2\ dx + \int_0^T \int_{{\bf R}^3} |\nabla u(t,x)|^2\ dx dt = \frac{1}{2} \int_{{\bf R}^3} |u(0,x)|^2\ dx.

This identity (and its consequences) provide essentially the only known a priori bound on solutions to the Navier-Stokes equations from large data and arbitrary times. Unfortunately, as discussed in this previous post, the quantities controlled by the energy identity are supercritical with respect to scaling, which is the fundamental obstacle that has defeated all attempts to solve the global regularity problem for Navier-Stokes without any additional assumptions on the data or solution (e.g. perturbative hypotheses, or a priori control on a critical norm such as the {L^\infty_t L^3_x} norm).

Our main result is then (slightly informally stated) as follows

Theorem 1 There exists an averaged version {\tilde B} of the bilinear operator {B}, of the form

\displaystyle  \tilde B(u,v) := \int_\Omega m_{3,\omega}(D) Rot_{3,\omega}

\displaystyle B( m_{1,\omega}(D) Rot_{1,\omega} u, m_{2,\omega}(D) Rot_{2,\omega} v )\ d\mu(\omega)

for some probability space {(\Omega, \mu)}, some spatial rotation operators {Rot_{i,\omega}} for {i=1,2,3}, and some Fourier multipliers {m_{i,\omega}} of order {0}, for which one still has the cancellation law

\displaystyle  \langle \tilde B(u,u), u \rangle = 0

and for which the averaged Navier-Stokes equation

\displaystyle  \partial_t u = \Delta u + \tilde B(u,u) \ \ \ \ \ (2)

admits solutions that blow up in finite time.

(There are some integrability conditions on the Fourier multipliers {m_{i,\omega}} required in the above theorem in order for the conclusion to be non-trivial, but I am omitting them here for sake of exposition.)

Because spatial rotations and Fourier multipliers of order {0} are bounded on most function spaces, {\tilde B} automatically obeys almost all of the upper bound estimates that {B} does. Thus, this theorem blocks any attempt to prove global regularity for the true Navier-Stokes equations which relies purely on the energy identity and on upper bound estimates for the nonlinearity; one must use some additional structure of the nonlinear operator {B} which is not shared by an averaged version {\tilde B}. Such additional structure certainly exists – for instance, the Navier-Stokes equation has a vorticity formulation involving only differential operators rather than pseudodifferential ones, whereas a general equation of the form (2) does not. However, “abstract” approaches to global regularity generally do not exploit such structure, and thus cannot be used to affirmatively answer the Navier-Stokes problem.

It turns out that the particular averaged bilinear operator {B} that we will use will be a finite linear combination of local cascade operators, which take the form

\displaystyle  C(u,v) := \sum_{n \in {\bf Z}} (1+\epsilon_0)^{5n/2} \langle u, \psi_{1,n} \rangle \langle v, \psi_{2,n} \rangle \psi_{3,n}

where {\epsilon_0>0} is a small parameter, {\psi_1,\psi_2,\psi_3} are Schwartz vector fields whose Fourier transform is supported on an annulus, and {\psi_{i,n}(x) := (1+\epsilon_0)^{3n/2} \psi_i( (1+\epsilon_0)^n x)} is an {L^2}-rescaled version of {\psi_i} (basically a “wavelet” of wavelength about {(1+\epsilon_0)^{-n}} centred at the origin). Such operators were essentially introduced by Katz and Pavlovic as dyadic models for {B}; they have the essentially the same scaling property as {B} (except that one can only scale along powers of {1+\epsilon_0}, rather than over all positive reals), and in fact they can be expressed as an average of {B} in the sense of the above theorem, as can be shown after a somewhat tedious amount of Fourier-analytic symbol manipulations.

If we consider nonlinearities {\tilde B} which are a finite linear combination of local cascade operators, then the equation (2) more or less collapses to a system of ODE in certain “wavelet coefficients” of {u}. The precise ODE that shows up depends on what precise combination of local cascade operators one is using. Katz and Pavlovic essentially considered a single cascade operator together with its “adjoint” (needed to preserve the energy identity), and arrived (more or less) at the system of ODE

\displaystyle  \partial_t X_n = - (1+\epsilon_0)^{2n} X_n + (1+\epsilon_0)^{\frac{5}{2}(n-1)} X_{n-1}^2 - (1+\epsilon_0)^{\frac{5}{2} n} X_n X_{n+1} \ \ \ \ \ (3)

where {X_n: [0,T] \rightarrow {\bf R}} are scalar fields for each integer {n}. (Actually, Katz-Pavlovic worked with a technical variant of this particular equation, but the differences are not so important for this current discussion.) Note that the quadratic terms on the RHS carry a higher exponent of {1+\epsilon_0} than the dissipation term; this reflects the supercritical nature of this evolution (the energy {\frac{1}{2} \sum_n X_n^2} is monotone decreasing in this flow, so the natural size of {X_n} given the control on the energy is {O(1)}). There is a slight technical issue with the dissipation if one wishes to embed (3) into an equation of the form (2), but it is minor and I will not discuss it further here.

In principle, if the {X_n} mode has size comparable to {1} at some time {t_n}, then energy should flow from {X_n} to {X_{n+1}} at a rate comparable to {(1+\epsilon_0)^{\frac{5}{2} n}}, so that by time {t_{n+1} \approx t_n + (1+\epsilon_0)^{-\frac{5}{2} n}} or so, most of the energy of {X_n} should have drained into the {X_{n+1}} mode (with hardly any energy dissipated). Since the series {\sum_{n \geq 1} (1+\epsilon_0)^{-\frac{5}{2} n}} is summable, this suggests finite time blowup for this ODE as the energy races ever more quickly to higher and higher modes. Such a scenario was indeed established by Katz and Pavlovic (and refined by Cheskidov) if the dissipation strength {(1+\epsilon)^{2n}} was weakened somewhat (the exponent {2} has to be lowered to be less than {\frac{5}{3}}). As mentioned above, this is enough to give a version of Theorem 1 in five and higher dimensions.

On the other hand, it was shown a few years ago by Barbato, Morandin, and Romito that (3) in fact admits global smooth solutions (at least in the dyadic case {\epsilon_0=1}, and assuming non-negative initial data). Roughly speaking, the problem is that as energy is being transferred from {X_n} to {X_{n+1}}, energy is also simultaneously being transferred from {X_{n+1}} to {X_{n+2}}, and as such the solution races off to higher modes a bit too prematurely, without absorbing all of the energy from lower modes. This weakens the strength of the blowup to the point where the moderately strong dissipation in (3) is enough to kill the high frequency cascade before a true singularity occurs. Because of this, the original Katz-Pavlovic model cannot quite be used to establish Theorem 1 in three dimensions. (Actually, the original Katz-Pavlovic model had some additional dispersive features which allowed for another proof of global smooth solutions, which is an unpublished result of Nazarov.)

To get around this, I had to “engineer” an ODE system with similar features to (3) (namely, a quadratic nonlinearity, a monotone total energy, and the indicated exponents of {(1+\epsilon_0)} for both the dissipation term and the quadratic terms), but for which the cascade of energy from scale {n} to scale {n+1} was not interrupted by the cascade of energy from scale {n+1} to scale {n+2}. To do this, I needed to insert a delay in the cascade process (so that after energy was dumped into scale {n}, it would take some time before the energy would start to transfer to scale {n+1}), but the process also needed to be abrupt (once the process of energy transfer started, it needed to conclude very quickly, before the delayed transfer for the next scale kicked in). It turned out that one could build a “quadratic circuit” out of some basic “quadratic gates” (analogous to how an electrical circuit could be built out of basic gates such as amplifiers or resistors) that achieved this task, leading to an ODE system essentially of the form

\displaystyle \partial_t X_{1,n} = - (1+\epsilon_0)^{2n} X_{1,n}

\displaystyle  + (1+\epsilon_0)^{5n/2} (- \epsilon^{-2} X_{3,n} X_{4,n} - \epsilon X_{1,n} X_{2,n} - \epsilon^2 \exp(-K^{10}) X_{1,n} X_{3,n}

\displaystyle  + K X_{4,n-1}^2)

\displaystyle  \partial_t X_{2,n} = - (1+\epsilon_0)^{2n} X_{2,n} + (1+\epsilon_0)^{5n/2} (\epsilon X_{1,n}^2 - \epsilon^{-1} K^{10} X_{3,n}^2)

\displaystyle  \partial_t X_{3,n} = - (1+\epsilon_0)^{2n} X_{3,n} + (1+\epsilon_0)^{5n/2} (\epsilon^2 \exp(-K^{10}) X_{1,n}^2

\displaystyle + \epsilon^{-1} K^{10} X_{2,n} X_{3,n} )

\displaystyle  \partial_t X_{4,n} =- (1+\epsilon_0)^{2n} X_{4,n} + (1+\epsilon_0)^{5n/2} (\epsilon^{-2} X_{3,n} X_{1,n}

\displaystyle - (1+\epsilon_0)^{5/2} K X_{4,n} X_{1,n+1})

where {K \geq 1} is a suitable large parameter and {\epsilon > 0} is a suitable small parameter (much smaller than {1/K}). To visualise the dynamics of such a system, I found it useful to describe this system graphically by a “circuit diagram” that is analogous (but not identical) to the circuit diagrams arising in electrical engineering:


The coupling constants here range widely from being very large to very small; in practice, this makes the {X_{2,n}} and {X_{3,n}} modes absorb very little energy, but exert a sizeable influence on the remaining modes. If a lot of energy is suddenly dumped into {X_{1,n}}, what happens next is roughly as follows: for a moderate period of time, nothing much happens other than a trickle of energy into {X_{2,n}}, which in turn causes a rapid exponential growth of {X_{3,n}} (from a very low base). After this delay, {X_{3,n}} suddenly crosses a certain threshold, at which point it causes {X_{1,n}} and {X_{4,n}} to exchange energy back and forth with extreme speed. The energy from {X_{4,n}} then rapidly drains into {X_{1,n+1}}, and the process begins again (with a slight loss in energy due to the dissipation). If one plots the total energy {E_n := \frac{1}{2} ( X_{1,n}^2 + X_{2,n}^2 + X_{3,n}^2 + X_{4,n}^2 )} as a function of time, it looks schematically like this:


As in the previous heuristic discussion, the time between cascades from one frequency scale to the next decay exponentially, leading to blowup at some finite time {T}. (One could describe the dynamics here as being similar to the famous “lighting the beacons” scene in the Lord of the Rings movies, except that (a) as each beacon gets ignited, the previous one is extinguished, as per the energy identity; (b) the time between beacon lightings decrease exponentially; and (c) there is no soundtrack.)

There is a real (but remote) possibility that this sort of construction can be adapted to the true Navier-Stokes equations. The basic blowup mechanism in the averaged equation is that of a von Neumann machine, or more precisely a construct (built within the laws of the inviscid evolution {\partial_t u = \tilde B(u,u)}) that, after some time delay, manages to suddenly create a replica of itself at a finer scale (and to largely erase its original instantiation in the process). In principle, such a von Neumann machine could also be built out of the laws of the inviscid form of the Navier-Stokes equations (i.e. the Euler equations). In physical terms, one would have to build the machine purely out of an ideal fluid (i.e. an inviscid incompressible fluid). If one could somehow create enough “logic gates” out of ideal fluid, one could presumably build a sort of “fluid computer”, at which point the task of building a von Neumann machine appears to reduce to a software engineering exercise rather than a PDE problem (providing that the gates are suitably stable with respect to perturbations, but (as with actual computers) this can presumably be done by converting the analog signals of fluid mechanics into a more error-resistant digital form). The key thing missing in this program (in both senses of the word) to establish blowup for Navier-Stokes is to construct the logic gates within the laws of ideal fluids. (Compare with the situation for cellular automata such as Conway’s “Game of Life“, in which Turing complete computers, universal constructors, and replicators have all been built within the laws of that game.)

I’ve just finished the first draft of my book “Expansion in finite simple groups of Lie type“, which is  based in the lecture notes for my graduate course on this topic that were previously posted on this blog.  It also contains some newer material, such as the notes on Lie algebras and Lie groups that I posted most recently here.

As always, corrections or comments are greatly appreciated (and errata will be collected at this page).

I’ve just uploaded to the arXiv my article “Algebraic combinatorial geometry: the polynomial method in arithmetic combinatorics, incidence combinatorics, and number theory“, submitted to the new journal “EMS surveys in the mathematical sciences“.  This is the first draft of a survey article on the polynomial method – a technique in combinatorics and number theory for controlling a relevant set of points by comparing it with the zero set of a suitably chosen polynomial, and then using tools from algebraic geometry (e.g. Bezout’s theorem) on that zero set. As such, the method combines algebraic geometry with combinatorial geometry, and could be viewed as the philosophy of a combined field which I dub “algebraic combinatorial geometry”.   There is also an important extension of this method when one is working overthe reals, in which methods from algebraic topology (e.g. the ham sandwich theorem and its generalisation to polynomials), and not just algebraic geometry, come into play also.

The polynomial method has been used independently many times in mathematics; for instance, it plays a key role in the proof of Baker’s theorem in transcendence theory, or Stepanov’s method in giving an elementary proof of the Riemann hypothesis for finite fields over curves; in combinatorics, the nullstellenatz of Alon is also another relatively early use of the polynomial method.  More recently, it underlies Dvir’s proof of the Kakeya conjecture over finite fields and Guth and Katz’s near-complete solution to the Erdos distance problem in the plane, and can be used to give a short proof of the Szemeredi-Trotter theorem.  One of the aims of this survey is to try to present all of these disparate applications of the polynomial method in a somewhat unified context; my hope is that there will eventually be a systematic foundation for algebraic combinatorial geometry which naturally contains all of these different instances the polynomial method (and also suggests new instances to explore); but the field is unfortunately not at that stage of maturity yet.

This is something of a first draft, so comments and suggestions are even more welcome than usual.  (For instance, I have already had my attention drawn to some additional uses of the polynomial method in the literature that I was not previously aware of.)

Emmanuel Breuillard, Ben Green, Bob Guralnick, and I have just uploaded to the arXiv our joint paper “Expansion in finite simple groups of Lie type“. This long-delayed paper (announced way back in 2010!) is a followup to our previous paper in which we showed that, with one possible exception, generic pairs of elements of a simple algebraic group (over an uncountable field) generated a free group which was strongly dense in the sense that any nonabelian subgroup of this group was Zariski dense. The main result of this paper is to establish the analogous result for finite simple groups of Lie type (as defined in the previous blog post) and bounded rank, namely that almost all pairs {a,b} of elements of such a group generate a Cayley graph which is a (two-sided) expander, with expansion constant bounded below by a quantity depending on the rank of the group. (Informally, this means that the random walk generated by {a,b} spreads out in logarithmic time to be essentially uniformly distributed across the group, as opposed for instance to being largely trapped in an algebraic subgroup. Thus if generic elements did not generate a strongly dense group, one would probably expect expansion to fail.)

There are also some related results established in the paper. Firstly, as we discovered after writing our first paper, there was one class of algebraic groups for which our demonstration of strongly dense subgroups broke down, namely the {Sp_4} groups in characteristic three. In the current paper we provide in a pair of appendices a new argument that covers this case (or more generally, {Sp_4} in odd characteristic), by first reducing to the case of affine groups {k^2 \rtimes SL_2(k)} (which can be found inside {Sp_4} as a subgroup) and then using a ping-pong argument (in a p-adic metric) in the latter context.

Secondly, we show that the distinction between one-sided expansion and two-sided expansion (see this set of lecture notes of mine for definitions) is erased in the context of Cayley graphs of bounded degree, in the sense that such graphs are one-sided expanders if and only if they are two-sided expanders (perhaps with slightly different expansion constants). The argument turns out to be an elementary combinatorial one, based on the “pivot” argument discussed in these lecture notes of mine.

Now to the main result of the paper, namely the expansion of random Cayley graphs. This result had previously been established for {SL_2} by Bourgain and Gamburd, and Ben, Emmanuel and I had used the Bourgain-Gamburd method to achieve the same result for Suzuki groups. For the other finite simple groups of Lie type, expander graphs had been constructed by Kassabov, Lubotzky, and Nikolov, but they required more than two generators, which were placed deterministically rather than randomly. (Here, I am skipping over a large number of other results on expanding Cayley graphs; see this survey of Lubotzsky for a fairly recent summary of developments.) The current paper also uses the “Bourgain-Gamburd machine”, as discussed in these lecture notes of mine, to demonstrate expansion. This machine shows how expansion of a Cayley graph follows from three basic ingredients, which we state informally as follows:

  • Non-concentration (A random walk in this graph does not concentrate in a proper subgroup);
  • Product theorem (A medium-sized subset of this group which is not trapped in a proper subgroup will expand under multiplication); and
  • Quasirandomness (The group has no small non-trivial linear representations).

Quasirandomness of arbitrary finite simple groups of Lie type was established many years ago (predating, in fact, the introduction of the term “quasirandomness” by Gowers for this property) by Landazuri-Seitz and Seitz-Zalesskii, and the product theorem was already established by Pyber-Szabo and independently by Breuillard, Green, and myself. So the main problem is to establish non-concentration: that for a random Cayley graph on a finite simple group {G} of Lie type, random walks did not concentrate in proper subgroups.

The first step was to classify the proper subgroups of {G}. Fortunately, these are all known; in particular, such groups are either contained in proper algebraic subgroups of the algebraic group containing {G} (or a bounded cover thereof) with bounded complexity, or are else arising (up to conjugacy) from a version {G(F')} of the same group {G =G(F)} associated to a proper subfield {F'} of the field {F} respectively; this follows for instance from the work of Larsen and Pink, but also can be deduced using the classification of finite simple groups, together with some work of Aschbacher, Liebeck-Seitz, and Nori. We refer to the two types of subgroups here as “structural subgroups” and “subfield subgroups”.

To preclude concentration in a structural subgroup, we use our previous result that generic elements of an algebraic group generate a strongly dense subgroup, and so do not concentrate in any algebraic subgroup. To translate this result from the algebraic group setting to the finite group setting, we need a Schwarz-Zippel lemma for finite simple groups of Lie type. This is straightforward for Chevalley groups, but turns out to be a bit trickier for the Steinberg and Suzuki-Ree groups, and we have to go back to the Chevalley-type parameterisation of such groups in terms of (twisted) one-parameter subgroups, that can be found for instance in the text of Carter; this “twisted Schwartz-Zippel lemma” may possibly have further application to analysis on twisted simple groups of Lie type. Unfortunately, the Schwartz-Zippel estimate becomes weaker in twisted settings, and particularly in the case of triality groups {{}^3 D_4(q)}, which require a somewhat ad hoc additional treatment that relies on passing to a simpler subgroup present in a triality group, namely a central product of two different {SL_2}‘s.

To rule out concentration in a conjugate of a subfield group, we repeat an argument we introduced in our Suzuki paper and pass to a matrix model and analyse the coefficients of the characteristic polynomial of words in this Cayley graph, to prevent them from concentrating in a subfield. (Note that these coefficients are conjugation-invariant.)

Van Vu and I have just uploaded to the arXiv our joint paper “Local universality of zeroes of random polynomials“. This paper is a sequel to our previous work on local universality of eigenvalues of (non-Hermitian) random matrices {M_n} with independent entries. One can re-interpret these previous results as a universality result for a certain type of random polynomial {f: {\bf C} \rightarrow {\bf C}}, namely the characteristic polynomial {f(z) = \hbox{det}(M_n-z)} of the random matrix {M_n}. In this paper, we consider the analogous problem for a different model of random polynomial, namely polynomials {f} with independent random coefficients. More precisely, we consider random polynomials {f = f_n} of the form

\displaystyle  f(z) = \sum_{i=0}^n c_i \xi_i z^n

where {c_0,\ldots,c_n \in {\bf C}} are deterministic complex coefficients, and {\xi_0,\ldots,\xi_n \equiv \xi} are independent identically distributed copies of a complex random variable {\xi}, which we normalise to have mean zero and variance one. For simplicity we will ignore the technical issue that the leading coefficient {c_n \xi_n} is allowed to vanish; then {f} has {n} zeroes {\zeta_1,\ldots,\zeta_n \in {\bf C}} (counting multiplicity), which can be viewed as a random point process {\Sigma = \{\zeta_1,\ldots,\zeta_n\}} in the complex plane. In analogy with other models (such as random matrix models), we expect the (suitably normalised) asymptotic statistics of this point process in the limit {n \rightarrow \infty} to be universal, in the sense that it is largely independent of the precise distribution of the atom variable {\xi}.

Our results are fairly general with regard to the choice of coefficients {c_i}, but we isolate three particular choice of coefficients that are particularly natural and well-studied in the literature:

  • Flat polynomials (or Weyl polynomials) in which {c_i := \frac{1}{\sqrt{i!}}}.
  • Elliptic polynomials (or binomial polynomials) in which {c_i := \sqrt{\binom{n}{i}}}.
  • Kac polynomials in which {c_i := 1}.

The flat and elliptic polynomials enjoy special symmetries in the model case when the atom distribution {\xi} is a complex Gaussian {N(0,1)_{\bf C}}. Indeed, the zeroes {\Sigma} of elliptic polynomials with complex Gaussian coefficients have a distribution which is invariant with respect to isometries {T: {\bf C} \cup \{\infty\} \rightarrow {\bf C} \cup \{\infty\}} of the Riemann sphere {{\bf C} \cup \{\infty\}} (thus {T\Sigma} has the same distribution as {\Sigma}), while the zeroes of the limiting case {\sum_{i=0}^\infty \frac{1}{\sqrt{i!}} \xi_i z^i} of the flat polynomials with complex Gaussian coefficients are similarly invariant with respect to isometries {T: {\bf C} \rightarrow {\bf C}} of the complex plane {{\bf C}}. (For a nice geometric proof of this facts, I recommend the nice book of Hough, Krishnapur, Peres, and Virag.)

The global (i.e. coarse-scale) distribution of zeroes of these polynomials is well understood, first in the case of gaussian distributions using the fundamental tool of the Kac-Rice formula, and then for more general choices of atom distribution in the recent work of Kabluchko and Zaporozhets. The zeroes of the flat polynomials are asymptotically distributed according to the circular law, normalised to be uniformly distributed on the disk {B(0,\sqrt{n})} of radius {\sqrt{n}} centred at the origin. To put it a bit informally, the zeroes are asymptotically distributed according to the measure {\frac{1}{\pi} 1_{|z| \leq \sqrt{n}} dz}, where {dz} denotes Lebesgue measure on the complex plane. One can non-rigorously see the scale {\sqrt{n}} appear by observing that when {|z|} is much larger than {\sqrt{n}}, we expect the leading term {\frac{1}{\sqrt{n!}} \xi_n z^n} of the flat polynomial {\sum_{i=0}^n \frac{1}{\sqrt{i!}} \xi_i z^i} to dominate, so that the polynomial should not have any zeroes in this region.

Similarly, the distribution of the elliptic polynomials is known to be asymptotically distributed according to a Cauchy-type distribution {\frac{1}{\pi} \frac{1}{1+|z|^2/n} dz}. The Kac polynomials {\sum_{i=0}^n \xi_i z^i} behave differently; the zeroes concentrate uniformly on the unit circle {|z|=1} (which is reasonable, given that one would expect the top order term {\xi_i z^i} to dominate for {|z| > 1} and the bottom order term {\xi_0} to dominate for {|z| < 1}). In particular, whereas the typical spacing between zeroes in the flat and elliptic cases would be expected to be comparable to {1}, the typical spacing between zeroes for a Kac polynomial would be expected instead to be comparable to {1/n}.

In our paper we studied the local distribution of zeroes at the scale of the typical spacing. In the case of polynomials with continuous complex atom disribution {\xi}, the natural statistic to measure here is the {k}-point correlation function {\rho^{(k)}(z_1,\ldots,z_k)}, which for distinct complex numbers {z_1,\ldots,z_k} can be defined as the probability that there is a zero in each of the balls {B(z_1,\varepsilon),\ldots,B(z_k,\varepsilon)} for some infinitesimal {\epsilon >0}, divided by the normalising factor {(\pi \epsilon^2)^k}. (One can also define a {k}-point correlation function in the case of a discrete distribution, but it will be a measure rather than a function in that case.) Our first main theorem is a general replacement principle which asserts, very roughly speaking, that the asymptotic {k}-point correlation functions of two random polynomials {f, \tilde f} will agree if the log-magnitudes {\log |f(z)|, \log |\tilde f(z)|} have asymptotically the same distribution (actually we have to consider the joint distribution of {\log |f(z_1)|, \ldots \log |f(z_k)|} for several points {z_1,\ldots,z_k}, but let us ignore this detail for now), and if the polynomials {f, \tilde f} obeys a “non-clustering property” which asserts, roughly speaking, that not too many of the zeroes of {f} can typically concentrate in a small ball. This replacement principle was implicit in our previous paper (and can be viewed as a local-scale version of the global-scale replacement principle in this earlier paper of ours). Specialising the replacement principle to the elliptic or flat polynomials, and using the Lindeberg swapping argument, we obtain a Two Moment Theorem that asserts, roughly speaking, that the asymptotic behaviour of the {k}-point correlation functions depends only on the first two moments of the real and imaginary components of {\xi}, as long as one avoids some regions of space where universality is expected to break down. (In particular, because {f(0) = c_0 \xi_0} does not have a universal distribution, one does not expect universality to hold near the origin; there is a similar problem near infinity.) Closely related results, by a slightly different method, have also been obtained recently by Ledoan, Merkli, and Starr. A similar result holds for the Kac polynomials after rescaling to account for the narrower spacing between zeroes.

We also have analogous results in the case of polynomials with real coefficients (so that the coefficients {c_i} and the atom distribution {\xi} are both real). In this case one expects to see a certain number of real zeroes, together with conjugate pairs of complex zeroes. Instead of the {k}-point correlation function {\rho^{(k)}(z_1,\ldots,z_k)}, the natural object of study is now the mixed {(k,l)}-point correlation function {\rho^{(k,l)}(x_1,\ldots,x_k,z_1,\ldots,z_l)} that (roughly speaking) controls how often one can find a real zero near the real numbers {x_1,\ldots,x_k}, and a complex zero near the points {z_1,\ldots,z_l}. It turns out that one can disentangle the real and strictly complex zeroes and obtain separate universality results for both zeroes, provided that at least one of the polynomials involved obeys a “weak repulsion estimate” that shows that the real zeroes do not cluster very close to each other (and that the complex zeroes do not cluster very close to their complex conjugates). Such an estimate is needed to avoid the situation of two nearby real zeroes “colliding” to create a (barely) complex zero and its complex conjugate, or the time reversal of such a collision. Fortunately, in the case of real gaussian polynomials one can use the Kac-Rice formula to establish such a weak repulsion estimate, allowing analogues of the above universality results for complex random polynomials in the real case. Among other things, this gives universality results for the number {N_{\bf R}} of real zeroes of a random flat or elliptic polynomial; it turns out this number is typically equal to {\frac{2}{\pi} \sqrt{n} + O(n^{1/2-c})} and {\frac{n} + O(n^{1/2-c})} respectively. (For Kac polynomials, the situation is somewhat different; it was already known that {N_{\bf R} = \frac{2}{\pi} \log n + o(\log n)} thanks to a long series of papers, starting with the foundational work of Kac and culminating in the work of Ibragimov and Maslova.)

While our methods are based on our earlier work on eigenvalues of random matrices, the situation with random polynomials is actually somewhat easier to deal with. This is because the log-magnitude {\log |f(z)|} of a random polynomial with independent coefficients is significantly easier to control than the log-determinant {\log |\hbox{det}(M-z)|} of a random matrix, as the former can be controlled by the central limit theorem, while the latter requires significantly more difficult arguments (in particular, bounds on the least singular value combined with Girko’s Hermitization trick). As such, one could view the current paper as an introduction to our more complicated previous paper, and with this in mind we have written the current paper to be self-contained (though this did make the paper somewhat lengthy, checking in at 68 pages).

The purpose of this post is to link to a short unpublished note of mine that I wrote back in 2010 but forgot to put on my web page at the time. Entitled “A physical space proof of the bilinear Strichartz and local smoothing estimates for the Schrodinger equation“, it gives a proof of two standard estimates for the free (linear) Schrodinger equation in flat Euclidean space, namely the bilinear Strichartz estimate and the local smoothing estimate, using primarily “physical space” methods such as integration by parts, instead of “frequency space” methods based on the Fourier transform, although a small amount of Fourier analysis (basically sectoral projection to make the Schrodinger waves move roughly in a given direction) is still needed.  This is somewhat in the spirit of an older paper of mine with Klainerman and Rodnianski doing something similar for the wave equation, and is also very similar to a paper of Planchon and Vega from 2009.  The hope was that by avoiding the finer properties of the Fourier transform, one could obtain a more robust argument which could also extend to nonlinear, non-free, or non-flat situations.   These notes were cited once or twice by some people that I had privately circulated them to, so I decided to put them online here for reference.

UPDATE, July 24: Fabrice Planchon has kindly supplied another note in which he gives a particularly simple proof of local smoothing in one dimension, and discusses some other variants of the method (related to the paper of Planchon and Vega cited earlier).

Tamar Ziegler and I have just uploaded to the arXiv our joint paper “A multi-dimensional Szemerédi theorem for the primes via a correspondence principle“. This paper is related to an earlier result of Ben Green and mine in which we established that the primes contain arbitrarily long arithmetic progressions. Actually, in that paper we proved a more general result:

Theorem 1 (Szemerédi’s theorem in the primes) Let {A} be a subset of the primes {{\mathcal P}} of positive relative density, thus {\limsup_{N \rightarrow \infty} \frac{|A \cap [N]|}{|{\mathcal P} \cap [N]|} > 0}. Then {A} contains arbitrarily long arithmetic progressions.

This result was based in part on an earlier paper of Green that handled the case of progressions of length three. With the primes replaced by the integers, this is of course the famous theorem of Szemerédi.

Szemerédi’s theorem has now been generalised in many different directions. One of these is the multidimensional Szemerédi theorem of Furstenberg and Katznelson, who used ergodic-theoretic techniques to show that any dense subset of {{\bf Z}^d} necessarily contained infinitely many constellations of any prescribed shape. Our main result is to relativise that theorem to the primes as well:

Theorem 2 (Multidimensional Szemerédi theorem in the primes) Let {d \geq 1}, and let {A} be a subset of the {d^{th}} Cartesian power {{\mathcal P}^d} of the primes of positive relative density, thus

\displaystyle  \limsup_{N \rightarrow \infty} \frac{|A \cap [N]^d|}{|{\mathcal P}^d \cap [N]^d|} > 0.

Then for any {v_1,\ldots,v_k \in {\bf Z}^d}, {A} contains infinitely many “constellations” of the form {a+r v_1, \ldots, a + rv_k} with {a \in {\bf Z}^k} and {r} a positive integer.

In the case when {A} is itself a Cartesian product of one-dimensional sets (in particular, if {A} is all of {{\mathcal P}^d}), this result already follows from Theorem 1, but there does not seem to be a similarly easy argument to deduce the general case of Theorem 2 from previous results. Simultaneously with this paper, an independent proof of Theorem 2 using a somewhat different method has been established by Cook, Maygar, and Titichetrakun.

The result is reminiscent of an earlier result of mine on finding constellations in the Gaussian primes (or dense subsets thereof). That paper followed closely the arguments of my original paper with Ben Green, namely it first enclosed (a W-tricked version of) the primes or Gaussian primes (in a sieve theoretic-sense) by a slightly larger set (or more precisely, a weight function {\nu}) of almost primes or almost Gaussian primes, which one could then verify (using methods closely related to the sieve-theoretic methods in the ongoing Polymath8 project) to obey certain pseudorandomness conditions, known as the linear forms condition and the correlation condition. Very roughly speaking, these conditions assert statements of the following form: if {n} is a randomly selected integer, then the events of {n+h_1,\ldots,n+h_k} simultaneously being an almost prime (or almost Gaussian prime) are approximately independent for most choices of {h_1,\ldots,h_k}. Once these conditions are satisfied, one can then run a transference argument (initially based on ergodic-theory methods, but nowadays there are simpler transference results based on the Hahn-Banach theorem, due to Gowers and Reingold-Trevisan-Tulsiani-Vadhan) to obtain relative Szemerédi-type theorems from their absolute counterparts.

However, when one tries to adapt these arguments to sets such as {{\mathcal P}^2}, a new difficulty occurs: the natural analogue of the almost primes would be the Cartesian square {{\mathcal A}^2} of the almost primes – pairs {(n,m)} whose entries are both almost primes. (Actually, for technical reasons, one does not work directly with a set of almost primes, but would instead work with a weight function such as {\nu(n) \nu(m)} that is concentrated on a set such as {{\mathcal A}^2}, but let me ignore this distinction for now.) However, this set {{\mathcal A}^2} does not enjoy as many pseudorandomness conditions as one would need for a direct application of the transference strategy to work. More specifically, given any fixed {h, k}, and random {(n,m)}, the four events

\displaystyle  (n,m) \in {\mathcal A}^2

\displaystyle  (n+h,m) \in {\mathcal A}^2

\displaystyle  (n,m+k) \in {\mathcal A}^2

\displaystyle  (n+h,m+k) \in {\mathcal A}^2

do not behave independently (as they would if {{\mathcal A}^2} were replaced for instance by the Gaussian almost primes), because any three of these events imply the fourth. This blocks the transference strategy for constellations which contain some right-angles to them (e.g. constellations of the form {(n,m), (n+r,m), (n,m+r)}) as such constellations soon turn into rectangles such as the one above after applying Cauchy-Schwarz a few times. (But a few years ago, Cook and Magyar showed that if one restricted attention to constellations which were in general position in the sense that any coordinate hyperplane contained at most one element in the constellation, then this obstruction does not occur and one can establish Theorem 2 in this case through the transference argument.) It’s worth noting that very recently, Conlon, Fox, and Zhao have succeeded in removing of the pseudorandomness conditions (namely the correlation condition) from the transference principle, leaving only the linear forms condition as the remaining pseudorandomness condition to be verified, but unfortunately this does not completely solve the above problem because the linear forms condition also fails for {{\mathcal A}^2} (or for weights concentrated on {{\mathcal A}^2}) when applied to rectangular patterns.

There are now two ways known to get around this problem and establish Theorem 2 in full generality. The approach of Cook, Magyar, and Titichetrakun proceeds by starting with one of the known proofs of the multidimensional Szemerédi theorem – namely, the proof that proceeds through hypergraph regularity and hypergraph removal – and attach pseudorandom weights directly within the proof itself, rather than trying to add the weights to the result of that proof through a transference argument. (A key technical issue is that weights have to be added to all the levels of the hypergraph – not just the vertices and top-order edges – in order to circumvent the failure of naive pseudorandomness.) As one has to modify the entire proof of the multidimensional Szemerédi theorem, rather than use that theorem as a black box, the Cook-Magyar-Titichetrakun argument is lengthier than ours; on the other hand, it is more general and does not rely on some difficult theorems about primes that are used in our paper.

In our approach, we continue to use the multidimensional Szemerédi theorem (or more precisely, the equivalent theorem of Furstenberg and Katznelson concerning multiple recurrence for commuting shifts) as a black box. The difference is that instead of using a transference principle to connect the relative multidimensional Szemerédi theorem we need to the multiple recurrence theorem, we instead proceed by a version of the Furstenberg correspondence principle, similar to the one that connects the absolute multidimensional Szemerédi theorem to the multiple recurrence theorem. I had discovered this approach many years ago in an unpublished note, but had abandoned it because it required an infinite number of linear forms conditions (in contrast to the transference technique, which only needed a finite number of linear forms conditions and (until the recent work of Conlon-Fox-Zhao) a correlation condition). The reason for this infinite number of conditions is that the correspondence principle has to build a probability measure on an entire {\sigma}-algebra; for this, it is not enough to specify the measure {\mu(A)} of a single set such as {A}, but one also has to specify the measure {\mu( T^{n_1} A \cap \ldots \cap T^{n_m} A)} of “cylinder sets” such as {T^{n_1} A \cap \ldots \cap T^{n_m} A} where {m} could be arbitrarily large. The larger {m} gets, the more linear forms conditions one needs to keep the correspondence under control.

With the sieve weights {\nu} we were using at the time, standard sieve theory methods could indeed provide a finite number of linear forms conditions, but not an infinite number, so my idea was abandoned. However, with my later work with Green and Ziegler on linear equations in primes (and related work on the Mobius-nilsequences conjecture and the inverse conjecture on the Gowers norm), Tamar and I realised that the primes themselves obey an infinite number of linear forms conditions, so one can basically use the primes (or a proxy for the primes, such as the von Mangoldt function {\Lambda}) as the enveloping sieve weight, rather than a classical sieve. Thus my old idea of using the Furstenberg correspondence principle to transfer Szemerédi-type theorems to the primes could actually be realised. In the one-dimensional case, this simply produces a much more complicated proof of Theorem 1 than the existing one; but it turns out that the argument works as well in higher dimensions and yields Theorem 2 relatively painlessly, except for the fact that it needs the results on linear equations in primes, the known proofs of which are extremely lengthy (and also require some of the transference machinery mentioned earlier). The problem of correlations in rectangles is avoided in the correspondence principle approach because one can compensate for such correlations by performing a suitable weighted limit to compute the measure {\mu( T^{n_1} A \cap \ldots \cap T^{n_m} A)} of cylinder sets, with each {m} requiring a different weighted correction. (This may be related to the Cook-Magyar-Titichetrakun strategy of weighting all of the facets of the hypergraph in order to recover pseudorandomness, although our contexts are rather different.)

Vitaly Bergelson, Tamar Ziegler, and I have just uploaded to the arXiv our joint paper “Multiple recurrence and convergence results associated to {{\bf F}_{p}^{\omega}}-actions“. This paper is primarily concerned with limit formulae in the theory of multiple recurrence in ergodic theory. Perhaps the most basic formula of this type is the mean ergodic theorem, which (among other things) asserts that if {(X,{\mathcal X}, \mu,T)} is a measure-preserving {{\bf Z}}-system (which, in this post, means that {(X,{\mathcal X}, \mu)} is a probability space and {T: X \mapsto X} is measure-preserving and invertible, thus giving an action {(T^n)_{n \in {\bf Z}}} of the integers), and {f,g \in L^2(X,{\mathcal X}, \mu)} are functions, and {X} is ergodic (which means that {L^2(X,{\mathcal X}, \mu)} contains no {T}-invariant functions other than the constants (up to almost everywhere equivalence, of course)), then the average

\displaystyle  \frac{1}{N} \sum_{n=1}^N \int_X f(x) g(T^n x)\ d\mu \ \ \ \ \ (1)

converges as {N \rightarrow \infty} to the expression

\displaystyle  (\int_X f(x)\ d\mu) (\int_X g(x)\ d\mu);

see e.g. this previous blog post. Informally, one can interpret this limit formula as an equidistribution result: if {x} is drawn at random from {X} (using the probability measure {\mu}), and {n} is drawn at random from {\{1,\ldots,N\}} for some large {N}, then the pair {(x, T^n x)} becomes uniformly distributed in the product space {X \times X} (using product measure {\mu \times \mu}) in the limit as {N \rightarrow \infty}.

If we allow {(X,\mu)} to be non-ergodic, then we still have a limit formula, but it is a bit more complicated. Let {{\mathcal X}^T} be the {T}-invariant measurable sets in {{\mathcal X}}; the {{\bf Z}}-system {(X, {\mathcal X}^T, \mu, T)} can then be viewed as a factor of the original system {(X, {\mathcal X}, \mu, T)}, which is equivalent (in the sense of measure-preserving systems) to a trivial system {(Z_0, {\mathcal Z}_0, \mu_{Z_0}, 1)} (known as the invariant factor) in which the shift is trivial. There is then a projection map {\pi_0: X \rightarrow Z_0} to the invariant factor which is a factor map, and the average (1) converges in the limit to the expression

\displaystyle  \int_{Z_0} (\pi_0)_* f(z) (\pi_0)_* g(z)\ d\mu_{Z_0}(x), \ \ \ \ \ (2)

where {(\pi_0)_*: L^2(X,{\mathcal X},\mu) \rightarrow L^2(Z_0,{\mathcal Z}_0,\mu_{Z_0})} is the pushforward map associated to the map {\pi_0: X \rightarrow Z_0}; see e.g. this previous blog post. We can interpret this as an equidistribution result. If {(x,T^n x)} is a pair as before, then we no longer expect complete equidistribution in {X \times X} in the non-ergodic, because there are now non-trivial constraints relating {x} with {T^n x}; indeed, for any {T}-invariant function {f: X \rightarrow {\bf C}}, we have the constraint {f(x) = f(T^n x)}; putting all these constraints together we see that {\pi_0(x) = \pi_0(T^n x)} (for almost every {x}, at least). The limit (2) can be viewed as an assertion that this constraint {\pi_0(x) = \pi_0(T^n x)} are in some sense the “only” constraints between {x} and {T^n x}, and that the pair {(x,T^n x)} is uniformly distributed relative to these constraints.

Limit formulae are known for multiple ergodic averages as well, although the statement becomes more complicated. For instance, consider the expression

\displaystyle  \frac{1}{N} \sum_{n=1}^N \int_X f(x) g(T^n x) h(T^{2n} x)\ d\mu \ \ \ \ \ (3)

for three functions {f,g,h \in L^\infty(X, {\mathcal X}, \mu)}; this is analogous to the combinatorial task of counting length three progressions in various sets. For simplicity we assume the system {(X,{\mathcal X},\mu,T)} to be ergodic. Naively one might expect this limit to then converge to

\displaystyle  (\int_X f\ d\mu) (\int_X g\ d\mu) (\int_X h\ d\mu)

which would roughly speaking correspond to an assertion that the triplet {(x,T^n x, T^{2n} x)} is asymptotically equidistributed in {X \times X \times X}. However, even in the ergodic case there can be additional constraints on this triplet that cannot be seen at the level of the individual pairs {(x,T^n x)}, {(x, T^{2n} x)}. The key obstruction here is that of eigenfunctions of the shift {T: X \rightarrow X}, that is to say non-trivial functions {f: X \rightarrow S^1} that obey the eigenfunction equation {Tf = \lambda f} almost everywhere for some constant (or {T}-invariant) {\lambda}. Each such eigenfunction generates a constraint

\displaystyle  f(x) \overline{f(T^n x)}^2 f(T^{2n} x) = 1 \ \ \ \ \ (4)

tying together {x}, {T^n x}, and {T^{2n} x}. However, it turns out that these are in some sense the only constraints on {x,T^n x, T^{2n} x} that are relevant for the limit (3). More precisely, if one sets {{\mathcal X}_1} to be the sub-algebra of {{\mathcal X}} generated by the eigenfunctions of {T}, then it turns out that the factor {(X, {\mathcal X}_1, \mu, T)} is isomorphic to a shift system {(Z_1, {\mathcal Z}_1, \mu_{Z_1}, x \mapsto x+\alpha)} known as the Kronecker factor, for some compact abelian group {Z_1 = (Z_1,+)} and some (irrational) shift {\alpha \in Z_1}; the factor map {\pi_1: X \rightarrow Z_1} pushes eigenfunctions forward to (affine) characters on {Z_1}. It is then known that the limit of (3) is

\displaystyle  \int_\Sigma (\pi_1)_* f(x_0) (\pi_1)_* g(x_1) (\pi_1)_* h(x_2)\ d\mu_\Sigma

where {\Sigma \subset Z_1^3} is the closed subgroup

\displaystyle  \Sigma = \{ (x_1,x_2,x_3) \in Z_1^3: x_1-2x_2+x_3=0 \}

and {\mu_\Sigma} is the Haar probability measure on {\Sigma}; see this previous blog post. The equation {x_1-2x_2+x_3=0} defining {\Sigma} corresponds to the constraint (4) mentioned earlier. Among other things, this limit formula implies Roth’s theorem, which in the context of ergodic theory is the assertion that the limit (or at least the limit inferior) of (3) is positive when {f=g=h} is non-negative and not identically vanishing.

If one considers a quadruple average

\displaystyle  \frac{1}{N} \sum_{n=1}^N \int_X f(x) g(T^n x) h(T^{2n} x) k(T^{3n} x)\ d\mu \ \ \ \ \ (5)

(analogous to counting length four progressions) then the situation becomes more complicated still, even in the ergodic case. In addition to the (linear) eigenfunctions that already showed up in the computation of the triple average (3), a new type of constraint also arises from quadratic eigenfunctions {f: X \rightarrow S^1}, which obey an eigenfunction equation {Tf = \lambda f} in which {\lambda} is no longer constant, but is now a linear eigenfunction. For such functions, {f(T^n x)} behaves quadratically in {n}, and one can compute the existence of a constraint

\displaystyle  f(x) \overline{f(T^n x)}^3 f(T^{2n} x)^3 \overline{f(T^{3n} x)} = 1 \ \ \ \ \ (6)

between {x}, {T^n x}, {T^{2n} x}, and {T^{3n} x} that is not detected at the triple average level. As it turns out, this is not the only type of constraint relevant for (5); there is a more general class of constraint involving two-step nilsystems which we will not detail here, but see e.g. this previous blog post for more discussion. Nevertheless there is still a similar limit formula to previous examples, involving a special factor {(Z_2, {\mathcal Z}_2, \mu_{Z_2}, S)} which turns out to be an inverse limit of two-step nilsystems; this limit theorem can be extracted from the structural theory in this paper of Host and Kra combined with a limit formula for nilsystems obtained by Lesigne, but will not be reproduced here. The pattern continues to higher averages (and higher step nilsystems); this was first done explicitly by Ziegler, and can also in principle be extracted from the structural theory of Host-Kra combined with nilsystem equidistribution results of Leibman. These sorts of limit formulae can lead to various recurrence results refining Roth’s theorem in various ways; see this paper of Bergelson, Host, and Kra for some examples of this.

The above discussion was concerned with {{\bf Z}}-systems, but one can adapt much of the theory to measure-preserving {G}-systems for other discrete countable abelian groups {G}, in which one now has a family {(T_g)_{g \in G}} of shifts indexed by {G} rather than a single shift, obeying the compatibility relation {T_{g+h}=T_g T_h}. The role of the intervals {\{1,\ldots,N\}} in this more general setting is replaced by that of Folner sequences. For arbitrary countable abelian {G}, the theory for double averages (1) and triple limits (3) is essentially identical to the {{\bf Z}}-system case. But when one turns to quadruple and higher limits, the situation becomes more complicated (and, for arbitrary {G}, still not fully understood). However one model case which is now well understood is the finite field case when {G = {\bf F}_p^\omega = \bigcup_{n=1}^\infty {\bf F}_p^n} is an infinite-dimensional vector space over a finite field {{\bf F}_p} (with the finite subspaces {{\bf F}_p^n} then being a good choice for the Folner sequence). Here, the analogue of the structural theory of Host and Kra was worked out by Vitaly, Tamar, and myself in these previous papers (treating the high characteristic and low characteristic cases respectively). In the finite field setting, it turns out that nilsystems no longer appear, and one only needs to deal with linear, quadratic, and higher order eigenfunctions (known collectively as phase polynomials). It is then natural to look for a limit formula that asserts, roughly speaking, that if {x} is drawn at random from a {{\bf F}_p^\omega}-system and {n} drawn randomly from a large subspace of {{\bf F}_p^\omega}, then the only constraints between {x, T^n x, \ldots, T^{(p-1)n} x} are those that arise from phase polynomials. The main theorem of this paper is to establish this limit formula (which, again, is a little complicated to state explicitly and will not be done here). In particular, we establish for the first time that the limit actually exists (a result which, for {{\bf Z}}-systems, was one of the main results of this paper of Host and Kra).

As a consequence, we can recover finite field analogues of most of the results of Bergelson-Host-Kra, though interestingly some of the counterexamples demonstrating sharpness of their results for {{\bf Z}}-systems (based on Behrend set constructions) do not seem to be present in the finite field setting (cf. this previous blog post on the cap set problem). In particular, we are able to largely settle the question of when one has a Khintchine-type theorem that asserts that for any measurable set {A} in an ergodic {{\bf F}_p^\omega}-system and any {\epsilon>0}, one has

\displaystyle  \mu( T_{c_1 n} A \cap \ldots \cap T_{c_k n} A ) > \mu(A)^k - \epsilon

for a syndetic set of {n}, where {c_1,\ldots,c_k \in {\bf F}_p} are distinct residue classes. It turns out that Khintchine-type theorems always hold for {k=1,2,3} (and for {k=1,2} ergodicity is not required), and for {k=4} it holds whenever {c_1,c_2,c_3,c_4} form a parallelogram, but not otherwise (though the counterexample here was such a painful computation that we ended up removing it from the paper, and may end up putting it online somewhere instead), and for larger {k} we could show that the Khintchine property failed for generic choices of {c_1,\ldots,c_k}, though the problem of determining exactly the tuples for which the Khintchine property failed looked to be rather messy and we did not completely settle it.

Emmanuel Breuillard, Ben Green, and I have just uploaded to the arXiv our survey “Small doubling in groups“, for the proceedings of the upcoming Erdos Centennial.  This is a short survey of the known results on classifying finite subsets A of an (abelian) additive group G = (G,+) or a (not necessarily abelian) multiplicative group G = (G,\cdot) that have small doubling in the sense that the sum set A+A or product set A \cdot A is small.  Such sets behave approximately like finite subgroups of G (and there is a closely related notion of an approximate group in which the analogy is even tighter) , and so this subject can be viewed as a sort of approximate version of finite group theory.  (Unfortunately, thus far the theory does not have much new to say about the classification of actual finite groups; progress has been largely made instead on classifying the (highly restricted) number of ways in which approximate groups can differ from a genuine group.)

In the classical case when G is the integers {\mathbb Z}, these sets were classified (in a qualitative sense, at least) by a celebrated theorem of Freiman, which roughly speaking says that such sets A are necessarily “commensurate” in some sense with a (generalised) arithmetic progression P of bounded rank.   There are a number of essentially equivalent ways to define what “commensurate” means here; for instance, in the original formulation of the theorem, one asks that A be a dense subset of P, but in modern formulations it is often more convenient to require instead that A be of comparable size to P and be covered by a bounded number of translates of P, or that A and P have an intersection that is of comparable size to both A and P (cf. the notion of commensurability in group theory).

Freiman’s original theorem was extended to more general abelian groups in a sequence of papers culminating in the paper of Green and Ruzsa that handled arbitrary abelian groups.   As such groups now contain non-trivial finite subgroups, the conclusion of the theorem must be  modified by allowing for “coset progressions” P+H, which can be viewed as “extensions”  of generalized arithmetic progressions P by genuine finite groups H.

The proof methods in these abelian results were Fourier-analytic in nature (except in the cases of sets of very small doubling, in which more combinatorial approaches can be applied, and there were also some geometric or combinatorial methods that gave some weaker structural results).  As such, it was a challenge to extend these results to nonabelian groups, although for various important special types of ambient group G (such as an linear group over a finite or infinite field) it turns out that one can use tools exploiting the special structure of those groups (e.g. for linear groups one would use tools from Lie theory and algebraic geometry) to obtain quite satisfactory results; see e.g. this survey of  Pyber and Szabo for the linear case.   When the ambient group G is completely arbitrary, it turns out the problem is closely related to the classical Hilbert’s fifth problem of determining the minimal requirements of a topological group in order for such groups to have Lie structure; this connection was first observed and exploited by Hrushovski, and then used by Breuillard, Green, and myself to obtain the analogue of Freiman’s theorem for an arbitrary nonabelian group.

This survey is too short to discuss in much detail the proof techniques used in these results (although the abelian case is discussed in this book of mine with Vu, and the nonabelian case discussed in this more recent book of mine), but instead focuses on the statements of the various known results, as well as some remaining open questions in the subject (in particular, there is substantial work left to be done in making the estimates more quantitative, particularly in the nonabelian setting).


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