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Asgar Jamneshan, Or Shalom, and myself have just uploaded to the arXiv our preprints “A Host–Kra {{\bf F}^\omega_2}-system of order 5 that is not Abramov of order 5, and non-measurability of the inverse theorem for the {U^6({\bf F}^n_2)} norm” and “The structure of totally disconnected Host–Kra–Ziegler factors, and the inverse theorem for the {U^k} Gowers uniformity norms on finite abelian groups of bounded torsion“. These two papers are both concerned with advancing the inverse theory for the Gowers norms and Gowers-Host-Kra seminorms; the first paper provides a counterexample in this theory (in particular disproving a conjecture of Bergelson, Ziegler and myself), and the second paper gives new positive results in the case when the underlying group is bounded torsion, or the ergodic system is totally disconnected. I discuss the two papers more below the fold.

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Tamar Ziegler and I have just uploaded to the arXiv our paper “Infinite partial sumsets in the primes“. This is a short paper inspired by a recent result of Kra, Moreira, Richter, and Robertson (discussed for instance in this Quanta article from last December) showing that for any set {A} of natural numbers of positive upper density, there exists a sequence {b_1 < b_2 < b_3 < \dots} of natural numbers and a shift {t} such that {b_i + b_j + t \in A} for all {i<j} this answers a question of Erdős). In view of the “transference principle“, it is then plausible to ask whether the same result holds if {A} is replaced by the primes. We can show the following results:

Theorem 1
  • (i) If the Hardy-Littlewood prime tuples conjecture (or the weaker conjecture of Dickson) is true, then there exists an increasing sequence {b_1 < b_2 < b_3 < \dots} of primes such that {b_i + b_j + 1} is prime for all {i < j}.
  • (ii) Unconditionally, there exist increasing sequences {a_1 < a_2 < \dots} and {b_1 < b_2 < \dots} of natural numbers such that {a_i + b_j} is prime for all {i<j}.
  • (iii) These conclusions fail if “prime” is replaced by “positive (relative) density subset of the primes” (even if the density is equal to 1).

We remark that it was shown by Balog that there (unconditionally) exist arbitrarily long but finite sequences {b_1 < \dots < b_k} of primes such that {b_i + b_j + 1} is prime for all {i < j \leq k}. (This result can also be recovered from the later results of Ben Green, myself, and Tamar Ziegler.) Also, it had previously been shown by Granville that on the Hardy-Littlewood prime tuples conjecture, there existed increasing sequences {a_1 < a_2 < \dots} and {b_1 < b_2 < \dots} of natural numbers such that {a_i+b_j} is prime for all {i,j}.

The conclusion of (i) is stronger than that of (ii) (which is of course consistent with the former being conditional and the latter unconditional). The conclusion (ii) also implies the well-known theorem of Maynard that for any given {k}, there exist infinitely many {k}-tuples of primes of bounded diameter, and indeed our proof of (ii) uses the same “Maynard sieve” that powers the proof of that theorem (though we use a formulation of that sieve closer to that in this blog post of mine). Indeed, the failure of (iii) basically arises from the failure of Maynard’s theorem for dense subsets of primes, simply by removing those clusters of primes that are unusually closely spaced.

Our proof of (i) was initially inspired by the topological dynamics methods used by Kra, Moreira, Richter, and Robertson, but we managed to condense it to a purely elementary argument (taking up only half a page) that makes no reference to topological dynamics and builds up the sequence {b_1 < b_2 < \dots} recursively by repeated application of the prime tuples conjecture.

The proof of (ii) takes up the majority of the paper. It is easiest to phrase the argument in terms of “prime-producing tuples” – tuples {(h_1,\dots,h_k)} for which there are infinitely many {n} with {n+h_1,\dots,n+h_k} all prime. Maynard’s theorem is equivalent to the existence of arbitrarily long prime-producing tuples; our theorem is equivalent to the stronger assertion that there exist an infinite sequence {h_1 < h_2 < \dots} such that every initial segment {(h_1,\dots,h_k)} is prime-producing. The main new tool for achieving this is the following cute measure-theoretic lemma of Bergelson:

Lemma 2 (Bergelson intersectivity lemma) Let {E_1,E_2,\dots} be subsets of a probability space {(X,\mu)} of measure uniformly bounded away from zero, thus {\inf_i \mu(E_i) > 0}. Then there exists a subsequence {E_{i_1}, E_{i_2}, \dots} such that

\displaystyle  \mu(E_{i_1} \cap \dots \cap E_{i_k} ) > 0

for all {k}.

This lemma has a short proof, though not an entirely obvious one. Firstly, by deleting a null set from {X}, one can assume that all finite intersections {E_{i_1} \cap \dots \cap E_{i_k}} are either positive measure or empty. Secondly, a routine application of Fatou’s lemma shows that the maximal function {\limsup_N \frac{1}{N} \sum_{i=1}^N 1_{E_i}} has a positive integral, hence must be positive at some point {x_0}. Thus there is a subsequence {E_{i_1}, E_{i_2}, \dots} whose finite intersections all contain {x_0}, thus have positive measure as desired by the previous reduction.

It turns out that one cannot quite combine the standard Maynard sieve with the intersectivity lemma because the events {E_i} that show up (which roughly correspond to the event that {n + h_i} is prime for some random number {n} (with a well-chosen probability distribution) and some shift {h_i}) have their probability going to zero, rather than being uniformly bounded from below. To get around this, we borrow an idea from a paper of Banks, Freiberg, and Maynard, and group the shifts {h_i} into various clusters {h_{i,1},\dots,h_{i,J_1}}, chosen in such a way that the probability that at least one of {n+h_{i,1},\dots,n+h_{i,J_1}} is prime is bounded uniformly from below. One then applies the Bergelson intersectivity lemma to those events and uses many applications of the pigeonhole principle to conclude.

Rachel Greenfeld and I have just uploaded to the arXiv our paper “A counterexample to the periodic tiling conjecture“. This is the full version of the result I announced on this blog a few months ago, in which we disprove the periodic tiling conjecture of Grünbaum-Shephard and Lagarias-Wang. The paper took a little longer than expected to finish, due to a technical issue that we did not realize at the time of the announcement that required a workaround.

In more detail: the original strategy, as described in the announcement, was to build a “tiling language” that was capable of encoding a certain “{p}-adic Sudoku puzzle”, and then show that the latter type of puzzle had only non-periodic solutions if {p} was a sufficiently large prime. As it turns out, the second half of this strategy worked out, but there was an issue in the first part: our tiling language was able (using {2}-group-valued functions) to encode arbitrary boolean relationships between boolean functions, and was also able (using {{\bf Z}/p{\bf Z}}-valued functions) to encode “clock” functions such as {n \mapsto n \hbox{ mod } p} that were part of our {p}-adic Sudoku puzzle, but we were not able to make these two types of functions “talk” to each other in the way that was needed to encode the {p}-adic Sudoku puzzle (the basic problem being that if {H} is a finite abelian {2}-group then there are no non-trivial subgroups of {H \times {\bf Z}/p{\bf Z}} that are not contained in {H} or trivial in the {{\bf Z}/p{\bf Z}} direction). As a consequence, we had to replace our “{p}-adic Sudoku puzzle” by a “{2}-adic Sudoku puzzle” which basically amounts to replacing the prime {p} by a sufficiently large power of {2} (we believe {2^{10}} will suffice). This solved the encoding issue, but the analysis of the {2}-adic Sudoku puzzles was a little bit more complicated than the {p}-adic case, for the following reason. The following is a nice exercise in analysis:

Theorem 1 (Linearity in three directions implies full linearity) Let {F: {\bf R}^2 \rightarrow {\bf R}} be a smooth function which is affine-linear on every horizontal line, diagonal (line of slope {1}), and anti-diagonal (line of slope {-1}). In other words, for any {c \in {\bf R}}, the functions {x \mapsto F(x,c)}, {x \mapsto F(x,c+x)}, and {x \mapsto F(x,c-x)} are each affine functions on {{\bf R}}. Then {F} is an affine function on {{\bf R}^2}.

Indeed, the property of being affine in three directions shows that the quadratic form associated to the Hessian {\nabla^2 F(x,y)} at any given point vanishes at {(1,0)}, {(1,1)}, and {(1,-1)}, and thus must vanish everywhere. In fact the smoothness hypothesis is not necessary; we leave this as an exercise to the interested reader. The same statement turns out to be true if one replaces {{\bf R}} with the cyclic group {{\bf Z}/p{\bf Z}} as long as {p} is odd; this is the key for us to showing that our {p}-adic Sudoku puzzles have an (approximate) two-dimensional affine structure, which on further analysis can then be used to show that it is in fact non-periodic. However, it turns out that the corresponding claim for cyclic groups {{\bf Z}/q{\bf Z}} can fail when {q} is a sufficiently large power of {2}! In fact the general form of functions {F: ({\bf Z}/q{\bf Z})^2 \rightarrow {\bf Z}/q{\bf Z}} that are affine on every horizontal line, diagonal, and anti-diagonal takes the form

\displaystyle  F(x,y) = Ax + By + C + D \frac{q}{4} y(x-y)

for some integer coefficients {A,B,C,D}. This additional “pseudo-affine” term {D \frac{q}{4} y(x-y)} causes some additional technical complications but ultimately turns out to be manageable.

During the writing process we also discovered that the encoding part of the proof becomes more modular and conceptual once one introduces two new definitions, that of an “expressible property” and a “weakly expressible property”. These concepts are somewhat analogous to that of {\Pi^0_0} sentences and {\Sigma^0_1} sentences in the arithmetic hierarchy, or to algebraic sets and semi-algebraic sets in real algebraic geometry. Roughly speaking, an expressible property is a property of a tuple of functions {f_w: G \rightarrow H_w}, {w \in {\mathcal W}} from an abelian group {G} to finite abelian groups {H_w}, such that the property can be expressed in terms of one or more tiling equations on the graph

\displaystyle  A := \{ (x, (f_w(x))_{w \in {\mathcal W}} \subset G \times \prod_{w \in {\mathcal W}} H_w.

For instance, the property that two functions {f,g: {\bf Z} \rightarrow H} differ by a constant can be expressed in terms of the tiling equation

\displaystyle  A \oplus (\{0\} \times H^2) = {\bf Z} \times H^2

(the vertical line test), as well as

\displaystyle  A \oplus (\{0\} \times \Delta \cup \{1\} \times (H^2 \backslash \Delta)) = G \times H^2,

where {\Delta = \{ (h,h): h \in H \}} is the diagonal subgroup of {H^2}. A weakly expressible property {P} is an existential quantification of some expressible property {P^*}, so that a tuple of functions {(f_w)_{w \in W}} obeys the property {P} if and only if there exists an extension of this tuple by some additional functions that obey the property {P^*}. It turns out that weakly expressible properties are closed under a number of useful operations, and allow us to easily construct quite complicated weakly expressible properties out of a “library” of simple weakly expressible properties, much as a complex computer program can be constructed out of simple library routines. In particular we will be able to “program” our Sudoku puzzle as a weakly expressible property.

Rachel Greenfeld and I have just uploaded to the arXiv our announcement “A counterexample to the periodic tiling conjecture“. This is an announcement of a longer paper that we are currently in the process of writing up (and hope to release in a few weeks), in which we disprove the periodic tiling conjecture of Grünbaum-Shephard and Lagarias-Wang. This conjecture can be formulated in both discrete and continuous settings:

Conjecture 1 (Discrete periodic tiling conjecture) Suppose that {F \subset {\bf Z}^d} is a finite set that tiles {{\bf Z}^d} by translations (i.e., {{\bf Z}^d} can be partitioned into translates of {F}). Then {F} also tiles {{\bf Z}^d} by translations periodically (i.e., the set of translations can be taken to be a periodic subset of {{\bf Z}^d}).

Conjecture 2 (Continuous periodic tiling conjecture) Suppose that {\Omega \subset {\bf R}^d} is a bounded measurable set of positive measure that tiles {{\bf R}^d} by translations up to null sets. Then {\Omega} also tiles {{\bf R}^d} by translations periodically up to null sets.

The discrete periodic tiling conjecture can be easily established for {d=1} by the pigeonhole principle (as first observed by Newman), and was proven for {d=2} by Bhattacharya (with a new proof given by Greenfeld and myself). The continuous periodic tiling conjecture was established for {d=1} by Lagarias and Wang. By an old observation of Hao Wang, one of the consequences of the (discrete) periodic tiling conjecture is that the problem of determining whether a given finite set {F \subset {\bf Z}^d} tiles by translations is (algorithmically and logically) decidable.

On the other hand, once one allows tilings by more than one tile, it is well known that aperiodic tile sets exist, even in dimension two – finite collections of discrete or continuous tiles that can tile the given domain by translations, but not periodically. Perhaps the most famous examples of such aperiodic tilings are the Penrose tilings, but there are many other constructions; for instance, there is a construction of Ammann, Grümbaum, and Shephard of eight tiles in {{\bf Z}^2} which tile aperiodically. Recently, Rachel and I constructed a pair of tiles in {{\bf Z}^d} that tiled a periodic subset of {{\bf Z}^d} aperiodically (in fact we could even make the tiling question logically undecidable in ZFC).

Our main result is then

Theorem 3 Both the discrete and continuous periodic tiling conjectures fail for sufficiently large {d}. Also, there is a finite abelian group {G_0} such that the analogue of the discrete periodic tiling conjecture for {{\bf Z}^2 \times G_0} is false.

This suggests that the techniques used to prove the discrete periodic conjecture in {{\bf Z}^2} are already close to the limit of their applicability, as they cannot handle even virtually two-dimensional discrete abelian groups such as {{\bf Z}^2 \times G_0}. The main difficulty is in constructing the counterexample in the {{\bf Z}^2 \times G_0} setting.

The approach starts by adapting some of the methods of a previous paper of Rachel and myself. The first step is make the problem easier to solve by disproving a “multiple periodic tiling conjecture” instead of the traditional periodic tiling conjecture. At present, Theorem 3 asserts the existence of a “tiling equation” {A \oplus F = {\bf Z}^2 \times G_0} (where one should think of {F} and {G_0} as given, and the tiling set {A} is known), which admits solutions, all of which are non-periodic. It turns out that it is enough to instead assert the existence of a system

\displaystyle  A \oplus F^{(m)} = {\bf Z}^2 \times G_0, m=1,\dots,M

of tiling equations, which admits solutions, all of which are non-periodic. This is basically because one can “stack” together a system of tiling equations into an essentially equivalent single tiling equation in a slightly larger group. The advantage of this reformulation is that it creates a “tiling language”, in which each sentence {A \oplus F^{(m)} = {\bf Z}^2 \times G_0} in the language expresses a different type of constraint on the unknown set {A}. The strategy then is to locate a non-periodic set {A} which one can try to “describe” by sentences in the tiling language that are obeyed by this non-periodic set, and which are “structured” enough that one can capture their non-periodic nature through enough of these sentences.

It is convenient to replace sets by functions, so that this tiling language can be translated to a more familiar language, namely the language of (certain types of) functional equations. The key point here is that the tiling equation

\displaystyle  A \oplus (\{0\} \times H) = G \times H

for some abelian groups {G, H} is precisely asserting that {A} is a graph

\displaystyle  A = \{ (x, f(x)): x \in G \}

of some function {f: G \rightarrow H} (this sometimes referred to as the “vertical line test” in U.S. undergraduate math classes). Using this translation, it is possible to encode a variety of functional equations relating one or more functions {f_i: G \rightarrow H} taking values in some finite group {H} (such as a cyclic group).

The non-periodic behaviour that we ended up trying to capture was that of a certain “{p}-adically structured function” {f_p: {\bf Z} \rightarrow ({\bf Z}/p{\bf Z})^\times} associated to some fixed and sufficiently large prime {p} (in fact for our arguments any prime larger than {48}, e.g., {p=53}, would suffice), defined by the formula

\displaystyle  f_p(n) := \frac{n}{p^{\nu_p(n)}} \hbox{ mod } p

for {n \neq 0} and {f_p(0)=1}, where {\nu_p(n)} is the number of times {p} divides {n}. In other words, {f_p(n)} is the last non-zero digit in the base {p} expansion of {n} (with the convention that the last non-zero digit of {0} is {1}). This function is not periodic, and yet obeys a lot of functional equations; for instance, one has {f_p(pn) = f_p(n)} for all {n}, and also {f_p(pn+j)=j} for {j=1,\dots,p-1} (and in fact these two equations, together with the condition {f_p(0)=1}, completely determine {f_p}). Here is what the function {f_p} looks like (for {p=5}):

It turns out that we cannot describe this one-dimensional non-periodic function directly via tiling equations. However, we can describe two-dimensional non-periodic functions such as {(n,m) \mapsto f_p(An+Bm+C)} for some coefficients {A,B,C} via a suitable system of tiling equations. A typical such function looks like this:

A feature of this function is that when one restricts to a row or diagonal of such a function, the resulting one-dimensional function exhibits “{p}-adic structure” in the sense that it behaves like a rescaled version of {f_p}; see the announcement for a precise version of this statement. It turns out that the converse is essentially true: after excluding some degenerate solutions in which the function is constant along one or more of the columns, all two-dimensional functions which exhibit {p}-adic structure along (non-vertical) lines must behave like one of the functions {(n,m) \mapsto f_p(An+Bm+C)} mentioned earlier, and in particular is non-periodic. The proof of this result is strongly reminiscent of the type of reasoning needed to solve a Sudoku puzzle, and so we have adopted some Sudoku-like terminology in our arguments to provide intuition and visuals. One key step is to perform a shear transformation to the puzzle so that many of the rows become constant, as displayed in this example,

and then perform a “Tetris” move of eliminating the constant rows to arrive at a secondary Sudoku puzzle which one then analyzes in turn:

It is the iteration of this procedure that ultimately generates the non-periodic {p}-adic structure.

Kaisa Matomäki, Xuancheng Shao, Joni Teräväinen, and myself have just uploaded to the arXiv our preprint “Higher uniformity of arithmetic functions in short intervals I. All intervals“. This paper investigates the higher order (Gowers) uniformity of standard arithmetic functions in analytic number theory (and specifically, the Möbius function {\mu}, the von Mangoldt function {\Lambda}, and the generalised divisor functions {d_k}) in short intervals {(X,X+H]}, where {X} is large and {H} lies in the range {X^{\theta+\varepsilon} \leq H \leq X^{1-\varepsilon}} for a fixed constant {0 < \theta < 1} (that one would like to be as small as possible). If we let {f} denote one of the functions {\mu, \Lambda, d_k}, then there is extensive literature on the estimation of short sums

\displaystyle  \sum_{X < n \leq X+H} f(n)

and some literature also on the estimation of exponential sums such as

\displaystyle  \sum_{X < n \leq X+H} f(n) e(-\alpha n)

for a real frequency {\alpha}, where {e(\theta) := e^{2\pi i \theta}}. For applications in the additive combinatorics of such functions {f}, it is also necessary to consider more general correlations, such as polynomial correlations

\displaystyle  \sum_{X < n \leq X+H} f(n) e(-P(n))

where {P: {\bf Z} \rightarrow {\bf R}} is a polynomial of some fixed degree, or more generally

\displaystyle  \sum_{X < n \leq X+H} f(n) \overline{F}(g(n) \Gamma)

where {G/\Gamma} is a nilmanifold of fixed degree and dimension (and with some control on structure constants), {g: {\bf Z} \rightarrow G} is a polynomial map, and {F: G/\Gamma \rightarrow {\bf C}} is a Lipschitz function (with some bound on the Lipschitz constant). Indeed, thanks to the inverse theorem for the Gowers uniformity norm, such correlations let one control the Gowers uniformity norm of {f} (possibly after subtracting off some renormalising factor) on such short intervals {(X,X+H]}, which can in turn be used to control other multilinear correlations involving such functions.

Traditionally, asymptotics for such sums are expressed in terms of a “main term” of some arithmetic nature, plus an error term that is estimated in magnitude. For instance, a sum such as {\sum_{X < n \leq X+H} \Lambda(n) e(-\alpha n)} would be approximated in terms of a main term that vanished (or is negligible) if {\alpha} is “minor arc”, but would be expressible in terms of something like a Ramanujan sum if {\alpha} was “major arc”, together with an error term. We found it convenient to cancel off such main terms by subtracting an approximant {f^\sharp} from each of the arithmetic functions {f} and then getting upper bounds on remainder correlations such as

\displaystyle  |\sum_{X < n \leq X+H} (f(n)-f^\sharp(n)) \overline{F}(g(n) \Gamma)| \ \ \ \ \ (1)

(actually for technical reasons we also allow the {n} variable to be restricted further to a subprogression of {(X,X+H]}, but let us ignore this minor extension for this discussion). There is some flexibility in how to choose these approximants, but we eventually found it convenient to use the following choices.

  • For the Möbius function {\mu}, we simply set {\mu^\sharp = 0}, as per the Möbius pseudorandomness conjecture. (One could choose a more sophisticated approximant in the presence of a Siegel zero, as I did with Joni in this recent paper, but we do not do so here.)
  • For the von Mangoldt function {\Lambda}, we eventually went with the Cramér-Granville approximant {\Lambda^\sharp(n) = \frac{W}{\phi(W)} 1_{(n,W)=1}}, where {W = \prod_{p < R} p} and {R = \exp(\log^{1/10} X)}.
  • For the divisor functions {d_k}, we used a somewhat complicated-looking approximant {d_k^\sharp(n) = \sum_{m \leq X^{\frac{k-1}{5k}}} P_m(\log n)} for some explicit polynomials {P_m}, chosen so that {d_k^\sharp} and {d_k} have almost exactly the same sums along arithmetic progressions (see the paper for details).

The objective is then to obtain bounds on sums such as (1) that improve upon the “trivial bound” that one can get with the triangle inequality and standard number theory bounds such as the Brun-Titchmarsh inequality. For {\mu} and {\Lambda}, the Siegel-Walfisz theorem suggests that it is reasonable to expect error terms that have “strongly logarithmic savings” in the sense that they gain a factor of {O_A(\log^{-A} X)} over the trivial bound for any {A>0}; for {d_k}, the Dirichlet hyperbola method suggests instead that one has “power savings” in that one should gain a factor of {X^{-c_k}} over the trivial bound for some {c_k>0}. In the case of the Möbius function {\mu}, there is an additional trick (introduced by Matomäki and Teräväinen) that allows one to lower the exponent {\theta} somewhat at the cost of only obtaining “weakly logarithmic savings” of shape {\log^{-c} X} for some small {c>0}.

Our main estimates on sums of the form (1) work in the following ranges:

  • For {\theta=5/8}, one can obtain strongly logarithmic savings on (1) for {f=\mu,\Lambda}, and power savings for {f=d_k}.
  • For {\theta=3/5}, one can obtain weakly logarithmic savings for {f = \mu, d_k}.
  • For {\theta=5/9}, one can obtain power savings for {f=d_3}.
  • For {\theta=1/3}, one can obtain power savings for {f=d_2}.

Conjecturally, one should be able to obtain power savings in all cases, and lower {\theta} down to zero, but the ranges of exponents and savings given here seem to be the limit of current methods unless one assumes additional hypotheses, such as GRH. The {\theta=5/8} result for correlation against Fourier phases {e(\alpha n)} was established previously by Zhan, and the {\theta=3/5} result for such phases and {f=\mu} was established previously by by Matomäki and Teräväinen.

By combining these results with tools from additive combinatorics, one can obtain a number of applications:

  • Direct insertion of our bounds in the recent work of Kanigowski, Lemanczyk, and Radziwill on the prime number theorem on dynamical systems that are analytic skew products gives some improvements in the exponents there.
  • We can obtain a “short interval” version of a multiple ergodic theorem along primes established by Frantzikinakis-Host-Kra and Wooley-Ziegler, in which we average over intervals of the form {(X,X+H]} rather than {[1,X]}.
  • We can obtain a “short interval” version of the “linear equations in primes” asymptotics obtained by Ben Green, Tamar Ziegler, and myself in this sequence of papers, where the variables in these equations lie in short intervals {(X,X+H]} rather than long intervals such as {[1,X]}.

We now briefly discuss some of the ingredients of proof of our main results. The first step is standard, using combinatorial decompositions (based on the Heath-Brown identity and (for the {\theta=3/5} result) the Ramaré identity) to decompose {\mu(n), \Lambda(n), d_k(n)} into more tractable sums of the following types:

  • Type {I} sums, which are basically of the form {\sum_{m \leq A:m|n} \alpha(m)} for some weights {\alpha(m)} of controlled size and some cutoff {A} that is not too large;
  • Type {II} sums, which are basically of the form {\sum_{A_- \leq m \leq A_+:m|n} \alpha(m)\beta(n/m)} for some weights {\alpha(m)}, {\beta(n)} of controlled size and some cutoffs {A_-, A_+} that are not too close to {1} or to {X};
  • Type {I_2} sums, which are basically of the form {\sum_{m \leq A:m|n} \alpha(m) d_2(n/m)} for some weights {\alpha(m)} of controlled size and some cutoff {A} that is not too large.

The precise ranges of the cutoffs {A, A_-, A_+} depend on the choice of {\theta}; our methods fail once these cutoffs pass a certain threshold, and this is the reason for the exponents {\theta} being what they are in our main results.

The Type {I} sums involving nilsequences can be treated by methods similar to those in this previous paper of Ben Green and myself; the main innovations are in the treatment of the Type {II} and Type {I_2} sums.

For the Type {II} sums, one can split into the “abelian” case in which (after some Fourier decomposition) the nilsequence {F(g(n)\Gamma)} is basically of the form {e(P(n))}, and the “non-abelian” case in which {G} is non-abelian and {F} exhibits non-trivial oscillation in a central direction. In the abelian case we can adapt arguments of Matomaki and Shao, which uses Cauchy-Schwarz and the equidistribution properties of polynomials to obtain good bounds unless {e(P(n))} is “major arc” in the sense that it resembles (or “pretends to be”) {\chi(n) n^{it}} for some Dirichlet character {\chi} and some frequency {t}, but in this case one can use classical multiplicative methods to control the correlation. It turns out that the non-abelian case can be treated similarly. After applying Cauchy-Schwarz, one ends up analyzing the equidistribution of the four-variable polynomial sequence

\displaystyle  (n,m,n',m') \mapsto (g(nm)\Gamma, g(n'm)\Gamma, g(nm') \Gamma, g(n'm'\Gamma))

as {n,m,n',m'} range in various dyadic intervals. Using the known multidimensional equidistribution theory of polynomial maps in nilmanifolds, one can eventually show in the non-abelian case that this sequence either has enough equidistribution to give cancellation, or else the nilsequence involved can be replaced with one from a lower dimensional nilmanifold, in which case one can apply an induction hypothesis.

For the type {I_2} sum, a model sum to study is

\displaystyle  \sum_{X < n \leq X+H} d_2(n) e(\alpha n)

which one can expand as

\displaystyle  \sum_{n,m: X < nm \leq X+H} e(\alpha nm).

We experimented with a number of ways to treat this type of sum (including automorphic form methods, or methods based on the Voronoi formula or van der Corput’s inequality), but somewhat to our surprise, the most efficient approach was an elementary one, in which one uses the Dirichlet approximation theorem to decompose the hyperbolic region {\{ (n,m) \in {\bf N}^2: X < nm \leq X+H \}} into a number of arithmetic progressions, and then uses equidistribution theory to establish cancellation of sequences such as {e(\alpha nm)} on the majority of these progressions. As it turns out, this strategy works well in the regime {H > X^{1/3+\varepsilon}} unless the nilsequence involved is “major arc”, but the latter case is treatable by existing methods as discussed previously; this is why the {\theta} exponent for our {d_2} result can be as low as {1/3}.

In a sequel to this paper (currently in preparation), we will obtain analogous results for almost all intervals {(x,x+H]} with {x} in the range {[X,2X]}, in which we will be able to lower {\theta} all the way to {0}.


Jan Grebik, Rachel Greenfeld, Vaclav Rozhon and I have just uploaded to the arXiv our preprint “Measurable tilings by abelian group actions“. This paper is related to an earlier paper of Rachel Greenfeld and myself concerning tilings of lattices {{\bf Z}^d}, but now we consider the more general situation of tiling a measure space {X} by a tile {A \subset X} shifted by a finite subset {F} of shifts of an abelian group {G = (G,+)} that acts in a measure-preserving (or at least quasi-measure-preserving) fashion on {X}. For instance, {X} could be a torus {{\bf T}^d = {\bf R}^d/{\bf Z}^d}, {A} could be a positive measure subset of that torus, and {G} could be the group {{\bf R}^d}, acting on {X} by translation.


If {F} is a finite subset of {G} with the property that the translates {f+A}, {f \in F} of {A \subset X} partition {X} up to null sets, we write {F \oplus A =_{a.e.} X}, and refer to this as a measurable tiling of {X} by {A} (with tiling set {F}). For instance, if {X} is the torus {{\bf T}^2}, we can create a measurable tiling with {A = [0,1/2]^2 \hbox{ mod } {\bf Z}^2} and {F = \{0,1/2\}^2}. Our main results are the following:

  • By modifying arguments from previous papers (including the one with Greenfeld mentioned above), we can establish the following “dilation lemma”: a measurable tiling {F \oplus A =_{a.e.} X} automatically implies further measurable tilings {rF \oplus A =_{a.e.} X}, whenever {r} is an integer coprime to all primes up to the cardinality {\# F} of {F}.
  • By averaging the above dilation lemma, we can also establish a “structure theorem” that decomposes the indicator function {1_A} of {A} into components, each of which are invariant with respect to a certain shift in {G}. We can establish this theorem in the case of measure-preserving actions on probability spaces via the ergodic theorem, but one can also generalize to other settings by using the device of “measurable medial means” (which relates to the concept of a universally measurable set).
  • By applying this structure theorem, we can show that all measurable tilings {F \oplus A = {\bf T}^1} of the one-dimensional torus {{\bf T}^1} are rational, in the sense that {F} lies in a coset of the rationals {{\bf Q} = {\bf Q}^1}. This answers a recent conjecture of Conley, Grebik, and Pikhurko; we also give an alternate proof of this conjecture using some previous results of Lagarias and Wang.
  • For tilings {F \oplus A = {\bf T}^d} of higher-dimensional tori, the tiling need not be rational. However, we can show that we can “slide” the tiling to be rational by giving each translate {f + A} of {A} a “velocity” {v_f \in {\bf R}^d}, and for every time {t}, the translates {f + tv_f + A} still form a partition of {{\bf T}^d} modulo null sets, and at time {t=1} the tiling becomes rational. In particular, if a set {A} can tile a torus in an irrational fashion, then it must also be able to tile the torus in a rational fashion.
  • In the two-dimensional case {d=2} one can arrange matters so that all the velocities {v_f} are parallel. If we furthermore assume that the tile {A} is connected, we can also show that the union of all the translates {f+A} with a common velocity {v_f = v} form a {v}-invariant subset of the torus.
  • Finally, we show that tilings {F \oplus A = {\bf Z}^d \times G} of a finitely generated discrete group {{\bf Z}^d \times G}, with {G} a finite group, cannot be constructed in a “local” fashion (we formalize this probabilistically using the notion of a “factor of iid process”) unless the tile {F} is contained in a single coset of {\{0\} \times G}. (Nonabelian local tilings, for instance of the sphere by rotations, are of interest due to connections with the Banach-Tarski paradox; see the aforementioned paper of Conley, Grebik, and Pikhurko. Unfortunately, our methods seem to break down completely in the nonabelian case.)

I’ve just uploaded to the arXiv my preprint “Perfectly packing a square by squares of nearly harmonic sidelength“. This paper concerns a variant of an old problem of Meir and Moser, who asks whether it is possible to perfectly pack squares of sidelength {1/n} for {n \geq 2} into a single square or rectangle of area {\sum_{n=2}^\infty \frac{1}{n^2} = \frac{\pi^2}{6} - 1}. (The following variant problem, also posed by Meir and Moser and discussed for instance in this MathOverflow post, is perhaps even more well known: is it possible to perfectly pack rectangles of dimensions {1/n \times 1/(n+1)} for {n \geq 1} into a single square of area {\sum_{n=1}^\infty \frac{1}{n(n+1)} = 1}?) For the purposes of this paper, rectangles and squares are understood to have sides parallel to the axes, and a packing is perfect if it partitions the region being packed up to sets of measure zero. As one partial result towards these problems, it was shown by Paulhus that squares of sidelength {1/n} for {n \geq 2} can be packed (not quite perfectly) into a single rectangle of area {\frac{\pi^2}{6} - 1 + \frac{1}{1244918662}}, and rectangles of dimensions {1/n \times 1/n+1} for {n \geq 1} can be packed (again not quite perfectly) into a single square of area {1 + \frac{1}{10^9+1}}. (Paulhus’s paper had some gaps in it, but these were subsequently repaired by Grzegorek and Januszewski.)

Another direction in which partial progress has been made is to consider instead the problem of packing squares of sidelength {n^{-t}}, {n \geq 1} perfectly into a square or rectangle of total area {\sum_{n=1}^\infty \frac{1}{n^{2t}}}, for some fixed constant {t > 1/2} (this lower bound is needed to make the total area {\sum_{n=1}^\infty \frac{1}{n^{2t}}} finite), with the aim being to get {t} as close to {1} as possible. Prior to this paper, the most recent advance in this direction was by Januszewski and Zielonka last year, who achieved such a packing in the range {1/2 < t \leq 2/3}.

In this paper we are able to get {t} arbitrarily close to {1} (which turns out to be a “critical” value of this parameter), but at the expense of deleting the first few tiles:

Theorem 1 If {1/2 < t < 1}, and {n_0} is sufficiently large depending on {t}, then one can pack squares of sidelength {n^{-t}}, {n \geq n_0} perfectly into a square of area {\sum_{n=n_0}^\infty \frac{1}{n^{2t}}}.

As in previous works, the general strategy is to execute a greedy algorithm, which can be described somewhat incompletely as follows.

  • Step 1: Suppose that one has already managed to perfectly pack a square {S} of area {\sum_{n=n_0}^\infty \frac{1}{n^{2t}}} by squares of sidelength {n^{-t}} for {n_0 \leq n < n_1}, together with a further finite collection {{\mathcal R}} of rectangles with disjoint interiors. (Initially, we would have {n_1=n_0} and {{\mathcal R} = \{S\}}, but these parameter will change over the course of the algorithm.)
  • Step 2: Amongst all the rectangles in {{\mathcal R}}, locate the rectangle {R} of the largest width (defined as the shorter of the two sidelengths of {R}).
  • Step 3: Pack (as efficiently as one can) squares of sidelength {n^{-t}} for {n_1 \leq n < n_2} into {R} for some {n_2>n_1}, and decompose the portion of {R} not covered by this packing into rectangles {{\mathcal R}'}.
  • Step 4: Replace {n_1} by {n_2}, replace {{\mathcal R}} by {({\mathcal R} \backslash \{R\}) \cup {\mathcal R}'}, and return to Step 1.

The main innovation of this paper is to perform Step 3 somewhat more efficiently than in previous papers.

The above algorithm can get stuck if one reaches a point where one has already packed squares of sidelength {1/n^t} for {n_0 \leq n < n_1}, but that all remaining rectangles {R} in {{\mathcal R}} have width less than {n_1^{-t}}, in which case there is no obvious way to fit in the next square. If we let {w(R)} and {h(R)} denote the width and height of these rectangles {R}, then the total area of the rectangles must be

\displaystyle  \sum_{R \in {\mathcal R}} w(R) h(R) = \sum_{n=n_0}^\infty \frac{1}{n^{2t}} - \sum_{n=n_0}^{n_1-1} \frac{1}{n^{2t}} \asymp n_1^{1-2t}

and the total perimeter {\mathrm{perim}({\mathcal R})} of these rectangles is

\displaystyle  \mathrm{perim}({\mathcal R}) = \sum_{R \in {\mathcal R}} 2(w(R)+h(R)) \asymp \sum_{R \in {\mathcal R}} h(R).

Thus we have

\displaystyle  n_1^{1-2t} \ll \mathrm{perim}({\mathcal R}) \sup_{R \in {\mathcal R}} w(R)

and so to ensure that there is at least one rectangle {R} with {w(R) \geq n_1^{-t}} it would be enough to have the perimeter bound

\displaystyle  \mathrm{perim}({\mathcal R}) \leq c n_1^{1-t}

for a sufficiently small constant {c>0}. It is here that we now see the critical nature of the exponent {t=1}: for {t<1}, the amount of perimeter we are permitted to have in the remaining rectangles increases as one progresses with the packing, but for {t=1} the amount of perimeter one is “budgeted” for stays constant (and for {t>1} the situation is even worse, in that the remaining rectangles {{\mathcal R}} should steadily decrease in total perimeter).

In comparison, the perimeter of the squares that one has already packed is equal to

\displaystyle  \sum_{n=n_0}^{n_1-1} 4 n^{-t}

which is comparable to {n_1^{1-t}} for {n_1} large (with the constants blowing up as {t} approaches the critical value of {1}). In previous algorithms, the total perimeter of the remainder rectangles {{\mathcal R}} was basically comparable to the perimeter of the squares already packed, and this is the main reason why the results only worked when {t} was sufficiently far away from {1}. In my paper, I am able to get the perimeter of {{\mathcal R}} significantly smaller than the perimeter of the squares already packed, by grouping those squares into lattice-like clusters (of about {M^2} squares arranged in an {M \times M} pattern), and sliding the squares in each cluster together to almost entirely eliminate the wasted space between each square, leaving only the space around the cluster as the main source of residual perimeter, which will be comparable to about {M n_1^{-t}} per cluster, as compared to the total perimeter of the squares in the cluster which is comparable to {M^2 n_1^{-t}}. This strategy is perhaps easiest to illustrate with a picture, in which {3 \times 4} squares {S_{i,j}} of slowly decreasing sidelength are packed together with relatively little wasted space:

By choosing the parameter {M} suitably large (and taking {n_0} sufficiently large depending on {M}), one can then prove the theorem. (In order to do some technical bookkeeping and to allow one to close an induction in the verification of the algorithm’s correctness, it is convenient to replace the perimeter {\sum_{R \in {\mathcal R}} 2(w(R)+h(R))} by a slightly weighted variant {\sum_{R \in {\mathcal R}} w(R)^\delta h(R)} for a small exponent {\delta}, but this is a somewhat artificial device that somewhat obscures the main ideas.)

Asgar Jamneshan and myself have just uploaded to the arXiv our preprint “The inverse theorem for the {U^3} Gowers uniformity norm on arbitrary finite abelian groups: Fourier-analytic and ergodic approaches“. This paper, which is a companion to another recent paper of ourselves and Or Shalom, studies the inverse theory for the third Gowers uniformity norm

\displaystyle  \| f \|_{U^3(G)}^8 = {\bf E}_{h_1,h_2,h_3,x \in G} \Delta_{h_1} \Delta_{h_2} \Delta_{h_3} f(x)

on an arbitrary finite abelian group {G}, where {\Delta_h f(x) := f(x+h) \overline{f(x)}} is the multiplicative derivative. Our main result is as follows:

Theorem 1 (Inverse theorem for {U^3(G)}) Let {G} be a finite abelian group, and let {f: G \rightarrow {\bf C}} be a {1}-bounded function with {\|f\|_{U^3(G)} \geq \eta} for some {0 < \eta \leq 1/2}. Then:
  • (i) (Correlation with locally quadratic phase) There exists a regular Bohr set {B(S,\rho) \subset G} with {|S| \ll \eta^{-O(1)}} and {\exp(-\eta^{-O(1)}) \ll \rho \leq 1/2}, a locally quadratic function {\phi: B(S,\rho) \rightarrow {\bf R}/{\bf Z}}, and a function {\xi: G \rightarrow \hat G} such that

    \displaystyle  {\bf E}_{x \in G} |{\bf E}_{h \in B(S,\rho)} f(x+h) e(-\phi(h)-\xi(x) \cdot h)| \gg \eta^{O(1)}.

  • (ii) (Correlation with nilsequence) There exists an explicit degree two filtered nilmanifold {H/\Lambda} of dimension {O(\eta^{-O(1)})}, a polynomial map {g: G \rightarrow H/\Lambda}, and a Lipschitz function {F: H/\Lambda \rightarrow {\bf C}} of constant {O(\exp(\eta^{-O(1)}))} such that

    \displaystyle  |{\bf E}_{x \in G} f(x) \overline{F}(g(x))| \gg \exp(-\eta^{-O(1)}).

Such a theorem was proven by Ben Green and myself in the case when {|G|} was odd, and by Samorodnitsky in the {2}-torsion case {G = {\bf F}_2^n}. In all cases one uses the “higher order Fourier analysis” techniques introduced by Gowers. After some now-standard manipulations (using for instance what is now known as the Balog-Szemerédi-Gowers lemma), one arrives (for arbitrary {G}) at an estimate that is roughly of the form

\displaystyle  |{\bf E}_{x \in G} {\bf E}_{h,k \in B(S,\rho)} f(x+h+k) b(x,k) b(x,h) e(-B(h,k))| \gg \eta^{O(1)}

where {b} denotes various {1}-bounded functions whose exact values are not too important, and {B: B(S,\rho) \times B(S,\rho) \rightarrow {\bf R}/{\bf Z}} is a symmetric locally bilinear form. The idea is then to “integrate” this form by expressing it in the form

\displaystyle  B(h,k) = \phi(h+k) - \phi(h) - \phi(k) \ \ \ \ \ (1)

for some locally quadratic {\phi: B(S,\rho) \rightarrow {\bf C}}; this then allows us to write the above correlation as

\displaystyle  |{\bf E}_{x \in G} {\bf E}_{h,k \in B(S,\rho)} f(x+h+k) e(-\phi(h+k)) b(x,k) b(x,h)| \gg \eta^{O(1)}

(after adjusting the {b} functions suitably), and one can now conclude part (i) of the above theorem using some linear Fourier analysis. Part (ii) follows by encoding locally quadratic phase functions as nilsequences; for this we adapt an algebraic construction of Manners.

So the key step is to obtain a representation of the form (1), possibly after shrinking the Bohr set {B(S,\rho)} a little if needed. This has been done in the literature in two ways:

  • When {|G|} is odd, one has the ability to divide by {2}, and on the set {2 \cdot B(S,\frac{\rho}{10}) = \{ 2x: x \in B(S,\frac{\rho}{10})\}} one can establish (1) with {\phi(h) := B(\frac{1}{2} h, h)}. (This is similar to how in single variable calculus the function {x \mapsto \frac{1}{2} x^2} is a function whose second derivative is equal to {1}.)
  • When {G = {\bf F}_2^n}, then after a change of basis one can take the Bohr set {B(S,\rho)} to be {{\bf F}_2^m} for some {m}, and the bilinear form can be written in coordinates as

    \displaystyle  B(h,k) = \sum_{1 \leq i,j \leq m} a_{ij} h_i k_j / 2 \hbox{ mod } 1

    for some {a_{ij} \in {\bf F}_2} with {a_{ij}=a_{ji}}. The diagonal terms {a_{ii}} cause a problem, but by subtracting off the rank one form {(\sum_{i=1}^m a_{ii} h_i) ((\sum_{i=1}^m a_{ii} k_i) / 2} one can write

    \displaystyle  B(h,k) = \sum_{1 \leq i,j \leq m} b_{ij} h_i k_j / 2 \hbox{ mod } 1

    on the orthogonal complement of {(a_{11},\dots,a_{mm})} for some coefficients {b_{ij}=b_{ji}} which now vanish on the diagonal: {b_{ii}=0}. One can now obtain (1) on this complement by taking

    \displaystyle  \phi(h) := \sum_{1 \leq i < j \leq m} b_{ij} h_i h_k / 2 \hbox{ mod } 1.

In our paper we can now treat the case of arbitrary finite abelian groups {G}, by means of the following two new ingredients:

  • (i) Using some geometry of numbers, we can lift the group {G} to a larger (possibly infinite, but still finitely generated) abelian group {G_S} with a projection map {\pi: G_S \rightarrow G}, and find a globally bilinear map {\tilde B: G_S \times G_S \rightarrow {\bf R}/{\bf Z}} on the latter group, such that one has a representation

    \displaystyle  B(\pi(x), \pi(y)) = \tilde B(x,y) \ \ \ \ \ (2)

    of the locally bilinear form {B} by the globally bilinear form {\tilde B} when {x,y} are close enough to the origin.
  • (ii) Using an explicit construction, one can show that every globally bilinear map {\tilde B: G_S \times G_S \rightarrow {\bf R}/{\bf Z}} has a representation of the form (1) for some globally quadratic function {\tilde \phi: G_S \rightarrow {\bf R}/{\bf Z}}.

To illustrate (i), consider the Bohr set {B(S,1/10) = \{ x \in {\bf Z}/N{\bf Z}: \|x/N\|_{{\bf R}/{\bf Z}} < 1/10\}} in {G = {\bf Z}/N{\bf Z}} (where {\|\|_{{\bf R}/{\bf Z}}} denotes the distance to the nearest integer), and consider a locally bilinear form {B: B(S,1/10) \times B(S,1/10) \rightarrow {\bf R}/{\bf Z}} of the form {B(x,y) = \alpha x y \hbox{ mod } 1} for some real number {\alpha} and all integers {x,y \in (-N/10,N/10)} (which we identify with elements of {G}. For generic {\alpha}, this form cannot be extended to a globally bilinear form on {G}; however if one lifts {G} to the finitely generated abelian group

\displaystyle  G_S := \{ (x,\theta) \in {\bf Z}/N{\bf Z} \times {\bf R}: \theta = x/N \hbox{ mod } 1 \}

(with projection map {\pi: (x,\theta) \mapsto x}) and introduces the globally bilinear form {\tilde B: G_S \times G_S \rightarrow {\bf R}/{\bf Z}} by the formula

\displaystyle  \tilde B((x,\theta),(y,\sigma)) = N^2 \alpha \theta \sigma \hbox{ mod } 1

then one has (2) when {\theta,\sigma} lie in the interval {(-1/10,1/10)}. A similar construction works for higher rank Bohr sets.

To illustrate (ii), the key case turns out to be when {G_S} is a cyclic group {{\bf Z}/N{\bf Z}}, in which case {\tilde B} will take the form

\displaystyle  \tilde B(x,y) = \frac{axy}{N} \hbox{ mod } 1

for some integer {a}. One can then check by direct construction that (1) will be obeyed with

\displaystyle  \tilde \phi(x) = \frac{a \binom{x}{2}}{N} - \frac{a x \binom{N}{2}}{N^2} \hbox{ mod } 1

regardless of whether {N} is even or odd. A variant of this construction also works for {{\bf Z}}, and the general case follows from a short calculation verifying that the claim (ii) for any two groups {G_S, G'_S} implies the corresponding claim (ii) for the product {G_S \times G'_S}.

This concludes the Fourier-analytic proof of Theorem 1. In this paper we also give an ergodic theory proof of (a qualitative version of) Theorem 1(ii), using a correspondence principle argument adapted from this previous paper of Ziegler, and myself. Basically, the idea is to randomly generate a dynamical system on the group {G}, by selecting an infinite number of random shifts {g_1, g_2, \dots \in G}, which induces an action of the infinitely generated free abelian group {{\bf Z}^\omega = \bigcup_{n=1}^\infty {\bf Z}^n} on {G} by the formula

\displaystyle  T^h x := x + \sum_{i=1}^\infty h_i g_i.

Much as the law of large numbers ensures the almost sure convergence of Monte Carlo integration, one can show that this action is almost surely ergodic (after passing to a suitable Furstenberg-type limit {X} where the size of {G} goes to infinity), and that the dynamical Host-Kra-Gowers seminorms of that system coincide with the combinatorial Gowers norms of the original functions. One is then well placed to apply an inverse theorem for the third Host-Kra-Gowers seminorm {U^3(X)} for {{\bf Z}^\omega}-actions, which was accomplished in the companion paper to this one. After doing so, one almost gets the desired conclusion of Theorem 1(ii), except that after undoing the application of the Furstenberg correspondence principle, the map {g: G \rightarrow H/\Lambda} is merely an almost polynomial rather than a polynomial, which roughly speaking means that instead of certain derivatives of {g} vanishing, they instead are merely very small outside of a small exceptional set. To conclude we need to invoke a “stability of polynomials” result, which at this level of generality was first established by Candela and Szegedy (though we also provide an independent proof here in an appendix), which roughly speaking asserts that every approximate polynomial is close in measure to an actual polynomial. (This general strategy is also employed in the Candela-Szegedy paper, though in the absence of the ergodic inverse theorem input that we rely upon here, the conclusion is weaker in that the filtered nilmanifold {H/\Lambda} is replaced with a general space known as a “CFR nilspace”.)

This transference principle approach seems to work well for the higher step cases (for instance, the stability of polynomials result is known in arbitrary degree); the main difficulty is to establish a suitable higher step inverse theorem in the ergodic theory setting, which we hope to do in future research.

Asgar Jamneshan, Or Shalom, and myself have just uploaded to the arXiv our preprint “The structure of arbitrary Conze–Lesigne systems“. As the title suggests, this paper is devoted to the structural classification of Conze-Lesigne systems, which are a type of measure-preserving system that are “quadratic” or of “complexity two” in a certain technical sense, and are of importance in the theory of multiple recurrence. There are multiple ways to define such systems; here is one. Take a countable abelian group {\Gamma} acting in a measure-preserving fashion on a probability space {(X,\mu)}, thus each group element {\gamma \in \Gamma} gives rise to a measure-preserving map {T^\gamma: X \rightarrow X}. Define the third Gowers-Host-Kra seminorm {\|f\|_{U^3(X)}} of a function {f \in L^\infty(X)} via the formula

\displaystyle  \|f\|_{U^3(X)}^8 := \lim_{n \rightarrow \infty} {\bf E}_{h_1,h_2,h_3 \in \Phi_n} \int_X \prod_{\omega_1,\omega_2,\omega_3 \in \{0,1\}}

\displaystyle {\mathcal C}^{\omega_1+\omega_2+\omega_3} f(T^{\omega_1 h_1 + \omega_2 h_2 + \omega_3 h_3} x)\ d\mu(x)

where {\Phi_n} is a Folner sequence for {\Gamma} and {{\mathcal C}: z \mapsto \overline{z}} is the complex conjugation map. One can show that this limit exists and is independent of the choice of Folner sequence, and that the {\| \|_{U^3(X)}} seminorm is indeed a seminorm. A Conze-Lesigne system is an ergodic measure-preserving system in which the {U^3(X)} seminorm is in fact a norm, thus {\|f\|_{U^3(X)}>0} whenever {f \in L^\infty(X)} is non-zero. Informally, this means that when one considers a generic parallelepiped in a Conze–Lesigne system {X}, the location of any vertex of that parallelepiped is more or less determined by the location of the other seven vertices. These are the important systems to understand in order to study “complexity two” patterns, such as arithmetic progressions of length four. While not all systems {X} are Conze-Lesigne systems, it turns out that they always have a maximal factor {Z^2(X)} that is a Conze-Lesigne system, known as the Conze-Lesigne factor or the second Host-Kra-Ziegler factor of the system, and this factor controls all the complexity two recurrence properties of the system.

The analogous theory in complexity one is well understood. Here, one replaces the {U^3(X)} norm by the {U^2(X)} norm

\displaystyle  \|f\|_{U^2(X)}^4 := \lim_{n \rightarrow \infty} {\bf E}_{h_1,h_2 \in \Phi_n} \int_X \prod_{\omega_1,\omega_2 \in \{0,1\}} {\mathcal C}^{\omega_1+\omega_2} f(T^{\omega_1 h_1 + \omega_2 h_2} x)\ d\mu(x)

and the ergodic systems for which {U^2} is a norm are called Kronecker systems. These systems are completely classified: a system is Kronecker if and only if it arises from a compact abelian group {Z} equipped with Haar probability measure and a translation action {T^\gamma \colon z \mapsto z + \phi(\gamma)} for some homomorphism {\phi: \Gamma \rightarrow Z} with dense image. Such systems can then be analyzed quite efficiently using the Fourier transform, and this can then be used to satisfactory analyze “complexity one” patterns, such as length three progressions, in arbitrary systems (or, when translated back to combinatorial settings, in arbitrary dense sets of abelian groups).

We return now to the complexity two setting. The most famous examples of Conze-Lesigne systems are (order two) nilsystems, in which the space {X} is a quotient {G/\Lambda} of a two-step nilpotent Lie group {G} by a lattice {\Lambda} (equipped with Haar probability measure), and the action is given by a translation {T^\gamma x = \phi(\gamma) x} for some group homomorphism {\phi: \Gamma \rightarrow G}. For instance, the Heisenberg {{\bf Z}}-nilsystem

\displaystyle  \begin{pmatrix} 1 & {\bf R} & {\bf R} \\ 0 & 1 & {\bf R} \\ 0 & 0 & 1 \end{pmatrix} / \begin{pmatrix} 1 & {\bf Z} & {\bf Z} \\ 0 & 1 & {\bf Z} \\ 0 & 0 & 1 \end{pmatrix}

with a shift of the form

\displaystyle  Tx = \begin{pmatrix} 1 & \alpha & 0 \\ 0 & 1 & \beta \\ 0 & 0 & 1 \end{pmatrix} x

for {\alpha,\beta} two real numbers with {1,\alpha,\beta} linearly independent over {{\bf Q}}, is a Conze-Lesigne system. As the base case of a well known result of Host and Kra, it is shown in fact that all Conze-Lesigne {{\bf Z}}-systems are inverse limits of nilsystems (previous results in this direction were obtained by Conze-Lesigne, Furstenberg-Weiss, and others). Similar results are known for {\Gamma}-systems when {\Gamma} is finitely generated, thanks to the thesis work of Griesmer (with further proofs by Gutman-Lian and Candela-Szegedy). However, this is not the case once {\Gamma} is not finitely generated; as a recent example of Shalom shows, Conze-Lesigne systems need not be the inverse limit of nilsystems in this case.

Our main result is that even in the infinitely generated case, Conze-Lesigne systems are still inverse limits of a slight generalisation of the nilsystem concept, in which {G} is a locally compact Polish group rather than a Lie group:

Theorem 1 (Classification of Conze-Lesigne systems) Let {\Gamma} be a countable abelian group, and {X} an ergodic measure-preserving {\Gamma}-system. Then {X} is a Conze-Lesigne system if and only if it is the inverse limit of translational systems {G/\Lambda}, where {G} is a nilpotent locally compact Polish group of nilpotency class two, and {\Lambda} is a lattice in {G} (and also a lattice in the commutator group {[G,G]}), with {G/\Lambda} equipped with the Haar probability measure and a translation action {T^\gamma x = \phi(\gamma) x} for some homomorphism {\phi: \Gamma \rightarrow G}.

In a forthcoming companion paper to this one, Asgar Jamneshan and I will use this theorem to derive an inverse theorem for the Gowers norm {U^3(G)} for an arbitrary finite abelian group {G} (with no restrictions on the order of {G}, in particular our result handles the case of even and odd {|G|} in a unified fashion). In principle, having a higher order version of this theorem will similarly allow us to derive inverse theorems for {U^{s+1}(G)} norms for arbitrary {s} and finite abelian {G}; we hope to investigate this further in future work.

We sketch some of the main ideas used to prove the theorem. The existing machinery developed by Conze-Lesigne, Furstenberg-Weiss, Host-Kra, and others allows one to describe an arbitrary Conze-Lesigne system as a group extension {Z \rtimes_\rho K}, where {Z} is a Kronecker system (a rotational system on a compact abelian group {Z = (Z,+)} and translation action {\phi: \Gamma \rightarrow Z}), {K = (K,+)} is another compact abelian group, and the cocycle {\rho = (\rho_\gamma)_{\gamma \in \Gamma}} is a collection of measurable maps {\rho_\gamma: Z \rightarrow K} obeying the cocycle equation

\displaystyle  \rho_{\gamma_1+\gamma_2}(x) = \rho_{\gamma_1}(T^{\gamma_2} x) + \rho_{\gamma_2}(x) \ \ \ \ \ (1)

for almost all {x \in Z}. Furthermore, {\rho} is of “type two”, which means in this concrete setting that it obeys an additional equation

\displaystyle  \rho_\gamma(x + z_1 + z_2) - \rho_\gamma(x+z_1) - \rho_\gamma(x+z_2) + \rho_\gamma(x) \ \ \ \ \ (2)

\displaystyle  = F(x + \phi(\gamma), z_1, z_2) - F(x,z_1,z_2)

for all {\gamma \in \Gamma} and almost all {x,z_1,z_2 \in Z}, and some measurable function {F: Z^3 \rightarrow K}; roughly speaking this asserts that {\phi_\gamma} is “linear up to coboundaries”. For technical reasons it is also convenient to reduce to the case where {Z} is separable. The problem is that the equation (2) is unwieldy to work with. In the model case when the target group {K} is a circle {{\bf T} = {\bf R}/{\bf Z}}, one can use some Fourier analysis to convert (2) into the more tractable Conze-Lesigne equation

\displaystyle  \rho_\gamma(x+z) - \rho_\gamma(x) = F_z(x+\phi(\gamma)) - F_z(x) + c_z(\gamma) \ \ \ \ \ (3)

for all {\gamma \in \Gamma}, all {z \in Z}, and almost all {x \in Z}, where for each {z}, {F_z: Z \rightarrow K} is a measurable function, and {c_z: \Gamma \rightarrow K} is a homomorphism. (For technical reasons it is often also convenient to enforce that {F_z, c_z} depend in a measurable fashion on {z}; this can always be achieved, at least when the Conze-Lesigne system is separable, but actually verifying that this is possible actually requires a certain amount of effort, which we devote an appendix to in our paper.) It is not difficult to see that (3) implies (2) for any group {K} (as long as one has the measurability in {z} mentioned previously), but the converse turns out to fail for some groups {K}, such as solenoid groups (e.g., inverse limits of {{\bf R}/2^n{\bf Z}} as {n \rightarrow \infty}), as was essentially shown by Rudolph. However, in our paper we were able to find a separate argument that also derived the Conze-Lesigne equation in the case of a cyclic group {K = \frac{1}{N}{\bf Z}/{\bf Z}}. Putting together the {K={\bf T}} and {K = \frac{1}{N}{\bf Z}/{\bf Z}} cases, one can then derive the Conze-Lesigne equation for arbitrary compact abelian Lie groups {K} (as such groups are isomorphic to direct products of finitely many tori and cyclic groups). As has been known for some time (see e.g., this paper of Host and Kra), once one has a Conze-Lesigne equation, one can more or less describe the system {X} as a translational system {G/\Lambda}, where the Host-Kra group {G} is the set of all pairs {(z, F_z)} that solve an equation of the form (3) (with these pairs acting on {X \equiv Z \rtimes_\rho K} by the law {(z,F_z) \cdot (x,k) := (x+z, k+F_z(x))}), and {\Lambda} is the stabiliser of a point in this system. This then establishes the theorem in the case when {K} is a Lie group, and the general case basically comes from the fact (from Fourier analysis or the Peter-Weyl theorem) that an arbitrary compact abelian group is an inverse limit of Lie groups. (There is a technical issue here in that one has to check that the space of translational system factors of {X} form a directed set in order to have a genuine inverse limit, but this can be dealt with by modifications of the tools mentioned here.)

There is an additional technical issue worth pointing out here (which unfortunately was glossed over in some previous work in the area). Because the cocycle equation (1) and the Conze-Lesigne equation (3) are only valid almost everywhere instead of everywhere, the action of {G} on {X} is technically only a near-action rather than a genuine action, and as such one cannot directly define {\Lambda} to be the stabiliser of a point without running into multiple problems. To fix this, one has to pass to a topological model of {X} in which the action becomes continuous, and the stabilizer becomes well defined, although one then has to work a little more to check that the action is still transitive. This can be done via Gelfand duality; we proceed using a mixture of a construction from this book of Host and Kra, and the machinery in this recent paper of Asgar and myself.

Now we discuss how to establish the Conze-Lesigne equation (3) in the cyclic group case {K = \frac{1}{N}{\bf Z}/{\bf Z}}. As this group embeds into the torus {{\bf T}}, it is easy to use existing methods obtain (3) but with the homomorphism {c_z} and the function {F_z} taking values in {{\bf R}/{\bf Z}} rather than in {\frac{1}{N}{\bf Z}/{\bf Z}}. The main task is then to fix up the homomorphism {c_z} so that it takes values in {\frac{1}{N}{\bf Z}/{\bf Z}}, that is to say that {Nc_z} vanishes. This only needs to be done locally near the origin, because the claim is easy when {z} lies in the dense subgroup {\phi(\Gamma)} of {Z}, and also because the claim can be shown to be additive in {z}. Near the origin one can leverage the Steinhaus lemma to make {c_z} depend linearly (or more precisely, homomorphically) on {z}, and because the cocycle {\rho} already takes values in {\frac{1}{N}{\bf Z}/{\bf Z}}, {N\rho} vanishes and {Nc_z} must be an eigenvalue of the system {Z}. But as {Z} was assumed to be separable, there are only countably many eigenvalues, and by another application of Steinhaus and linearity one can then make {Nc_z} vanish on an open neighborhood of the identity, giving the claim.

Joni Teräväinen and I have just uploaded to the arXiv our preprint “The Hardy–Littlewood–Chowla conjecture in the presence of a Siegel zero“. This paper is a development of the theme that certain conjectures in analytic number theory become easier if one makes the hypothesis that Siegel zeroes exist; this places one in a presumably “illusory” universe, since the widely believed Generalised Riemann Hypothesis (GRH) precludes the existence of such zeroes, yet this illusory universe seems remarkably self-consistent and notoriously impossible to eliminate from one’s analysis.

For the purposes of this paper, a Siegel zero is a zero {\beta} of a Dirichlet {L}-function {L(\cdot,\chi)} corresponding to a primitive quadratic character {\chi} of some conductor {q_\chi}, which is close to {1} in the sense that

\displaystyle  \beta = 1 - \frac{1}{\eta \log q_\chi}

for some large {\eta \gg 1} (which we will call the quality) of the Siegel zero. The significance of these zeroes is that they force the Möbius function {\mu} and the Liouville function {\lambda} to “pretend” to be like the exceptional character {\chi} for primes of magnitude comparable to {q_\chi}. Indeed, if we define an exceptional prime to be a prime {p^*} in which {\chi(p^*) \neq -1}, then very few primes near {q_\chi} will be exceptional; in our paper we use some elementary arguments to establish the bounds

\displaystyle  \sum_{q_\chi^{1/2+\varepsilon} < p^* \leq x} \frac{1}{p^*} \ll_\varepsilon \frac{\log x}{\eta \log q_\chi} \ \ \ \ \ (1)

for any {x \geq q_\chi^{1/2+\varepsilon}} and {\varepsilon>0}, where the sum is over exceptional primes in the indicated range {q_\chi^{1/2+\varepsilon} < p^* \leq x}; this bound is non-trivial for {x} as large as {q_\chi^{\eta^{1-\varepsilon}}}. (See Section 1 of this blog post for some variants of this argument, which were inspired by work of Heath-Brown.) There is also a companion bound (somewhat weaker) that covers a range of {p^*} a little bit below {q_\chi^{1/2}}.

One of the early influential results in this area was the following result of Heath-Brown, which I previously blogged about here:

Theorem 1 (Hardy-Littlewood assuming Siegel zero) Let {h} be a fixed natural number. Suppose one has a Siegel zero {\beta} associated to some conductor {q_\chi}. Then we have

\displaystyle  \sum_{n \leq x} \Lambda(n) \Lambda(n+h) = ({\mathfrak S} + O( \frac{1}{\log\log \eta} )) x

for all {q_\chi^{250} \leq x \leq q_\chi^{300}}, where {\Lambda} is the von Mangoldt function and {{\mathfrak S}} is the singular series

\displaystyle  {\mathfrak S} = \prod_{p|h} \frac{p}{p-1} \prod_{p \nmid h} (1 - \frac{1}{(p-1)^2})

In particular, Heath-Brown showed that if there are infinitely many Siegel zeroes, then there are also infinitely many twin primes, with the correct asymptotic predicted by the Hardy-Littlewood prime tuple conjecture at infinitely many scales.

Very recently, Chinis established an analogous result for the Chowla conjecture (building upon earlier work of Germán and Katai):

Theorem 2 (Chowla assuming Siegel zero) Let {h_1,\dots,h_\ell} be distinct fixed natural numbers. Suppose one has a Siegel zero {\beta} associated to some conductor {q_\chi}. Then one has

\displaystyle  \sum_{n \leq x} \lambda(n+h_1) \dots \lambda(n+h_\ell) \ll \frac{x}{(\log\log \eta)^{1/2} (\log \eta)^{1/12}}

in the range {q_\chi^{10} \leq x \leq q_\chi^{\log\log \eta / 3}}, where {\lambda} is the Liouville function.

In our paper we unify these results and also improve the quantitative estimates and range of {x}:

Theorem 3 (Hardy-Littlewood-Chowla assuming Siegel zero) Let {h_1,\dots,h_k,h'_1,\dots,h'_\ell} be distinct fixed natural numbers with {k \leq 2}. Suppose one has a Siegel zero {\beta} associated to some conductor {q_\chi}. Then one has

\displaystyle  \sum_{n \leq x} \Lambda(n+h_1) \dots \Lambda(n+h_k) \lambda(n+h'_1) \dots \lambda(n+h'_\ell)

\displaystyle = ({\mathfrak S} + O_\varepsilon( \frac{1}{\log^{1/10\max(1,k)} \eta} )) x

for

\displaystyle  q_\chi^{10k+\frac{1}{2}+\varepsilon} \leq x \leq q_\chi^{\eta^{1/2}}

for any fixed {\varepsilon>0}.

Our argument proceeds by a series of steps in which we replace {\Lambda} and {\lambda} by more complicated looking, but also more tractable, approximations, until the correlation is one that can be computed in a tedious but straightforward fashion by known techniques. More precisely, the steps are as follows:

  • (i) Replace the Liouville function {\lambda} with an approximant {\lambda_{\mathrm{Siegel}}}, which is a completely multiplicative function that agrees with {\lambda} at small primes and agrees with {\chi} at large primes.
  • (ii) Replace the von Mangoldt function {\Lambda} with an approximant {\Lambda_{\mathrm{Siegel}}}, which is the Dirichlet convolution {\chi * \log} multiplied by a Selberg sieve weight {\nu} to essentially restrict that convolution to almost primes.
  • (iii) Replace {\lambda_{\mathrm{Siegel}}} with a more complicated truncation {\lambda_{\mathrm{Siegel}}^\sharp} which has the structure of a “Type I sum”, and which agrees with {\lambda_{\mathrm{Siegel}}} on numbers that have a “typical” factorization.
  • (iv) Replace the approximant {\Lambda_{\mathrm{Siegel}}} with a more complicated approximant {\Lambda_{\mathrm{Siegel}}^\sharp} which has the structure of a “Type I sum”.
  • (v) Now that all terms in the correlation have been replaced with tractable Type I sums, use standard Euler product calculations and Fourier analysis, similar in spirit to the proof of the pseudorandomness of the Selberg sieve majorant for the primes in this paper of Ben Green and myself, to evaluate the correlation to high accuracy.

Steps (i), (ii) proceed mainly through estimates such as (1) and standard sieve theory bounds. Step (iii) is based primarily on estimates on the number of smooth numbers of a certain size.

The restriction {k \leq 2} in our main theorem is needed only to execute step (iv) of this step. Roughly speaking, the Siegel approximant {\Lambda_{\mathrm{Siegel}}} to {\Lambda} is a twisted, sieved version of the divisor function {\tau}, and the types of correlation one is faced with at the start of step (iv) are a more complicated version of the divisor correlation sum

\displaystyle  \sum_{n \leq x} \tau(n+h_1) \dots \tau(n+h_k).

For {k=1} this sum can be easily controlled by the Dirichlet hyperbola method. For {k=2} one needs the fact that {\tau} has a level of distribution greater than {1/2}; in fact Kloosterman sum bounds give a level of distribution of {2/3}, a folklore fact that seems to have first been observed by Linnik and Selberg. We use a (notationally more complicated) version of this argument to treat the sums arising in (iv) for {k \leq 2}. Unfortunately for {k > 2} there are no known techniques to unconditionally obtain asymptotics, even for the model sum

\displaystyle  \sum_{n \leq x} \tau(n) \tau(n+1) \tau(n+2),

although we do at least have fairly convincing conjectures as to what the asymptotics should be. Because of this, it seems unlikely that one will be able to relax the {k \leq 2} hypothesis in our main theorem at our current level of understanding of analytic number theory.

Step (v) is a tedious but straightforward sieve theoretic computation, similar in many ways to the correlation estimates of Goldston and Yildirim used in their work on small gaps between primes (as discussed for instance here), and then also used by Ben Green and myself to locate arithmetic progressions in primes.

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