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Homogenization of iterated singular integrals with applications to random quasiconformal maps

22 June, 2020 in math.CA, math.CV, paper | Tags: Beltrami equation, Eero Saksman, Kari Astala, quasiconformal mapping, singular integrals, Steffen Rohde | by Terence Tao | 10 comments

Kari Astala, Steffen Rohde, Eero Saksman and I have (finally!) uploaded to the arXiv our preprint “Homogenization of iterated singular integrals with applications to random quasiconformal maps“. This project started (and was largely completed) over a decade ago, but for various reasons it was not finalised until very recently. The motivation for this project was to study the behaviour of “random” quasiconformal maps. Recall that a (smooth) quasiconformal map is a homeomorphism ${f: {\bf C} \rightarrow {\bf C}}$ that obeys the Beltrami equation

$\displaystyle \frac{\partial f}{\partial \overline{z}} = \mu \frac{\partial f}{\partial z}$

for some Beltrami coefficient ${\mu: {\bf C} \rightarrow D(0,1)}$ ; this can be viewed as a deformation of the Cauchy-Riemann equation ${\frac{\partial f}{\partial \overline{z}} = 0}$ . Assuming that ${f(z)}$ is asymptotic to ${z}$ at infinity, one can (formally, at least) solve for ${f}$ in terms of ${\mu}$ using the Beurling transform

$\displaystyle Tf(z) := \frac{\partial}{\partial z}(\frac{\partial f}{\partial \overline{z}})^{-1}(z) = -\frac{1}{\pi} p.v. \int_{\bf C} \frac{f(w)}{(w-z)^2}\ dw$

by the Neumann series

$\displaystyle \frac{\partial f}{\partial \overline{z}} = \mu + \mu T \mu + \mu T \mu T \mu + \dots.$

We looked at the question of the asymptotic behaviour of ${f}$ if ${\mu = \mu_\delta}$ is a random field that oscillates at some fine spatial scale ${\delta>0}$ . A simple model to keep in mind is

$\displaystyle \mu_\delta(z) = \varphi(z) \sum_{n \in {\bf Z}^2} \epsilon_n 1_{n\delta + [0,\delta]^2}(z) \ \ \ \ \ (1)$

where ${\epsilon_n = \pm 1}$ are independent random signs and ${\varphi: {\bf C} \rightarrow D(0,1)}$ is a bump function. For models such as these, we show that a homogenisation occurs in the limit ${\delta \rightarrow 0}$ ; each multilinear expression

$\displaystyle \mu_\delta T \mu_\delta \dots T \mu_\delta \ \ \ \ \ (2)$

converges weakly in probability (and almost surely, if we restrict ${\delta}$ to a lacunary sequence) to a deterministic limit, and the associated quasiconformal map ${f = f_\delta}$ similarly converges weakly in probability (or almost surely). (Results of this latter type were also recently obtained by Ivrii and Markovic by a more geometric method which is simpler, but is applied to a narrower class of Beltrami coefficients.) In the specific case (1), the limiting quasiconformal map is just the identity map ${f(z)=z}$ , but if for instance replaces the ${\epsilon_n}$ by non-symmetric random variables then one can have significantly more complicated limits. The convergence theorem for multilinear expressions such as is not specific to the Beurling transform ${T}$ ; any other translation and dilation invariant singular integral can be used here.

The random expression (2) is somewhat reminiscent of a moment of a random matrix, and one can start computing it analogously. For instance, if one has a decomposition ${\mu_\delta = \sum_{n \in {\bf Z}^2} \mu_{\delta,n}}$ such as (1), then (2) expands out as a sum

$\displaystyle \sum_{n_1,\dots,n_k \in {\bf Z}^2} \mu_{\delta,n_1} T \mu_{\delta,n_2} \dots T \mu_{\delta,n_k}$

The random fluctuations of this sum can be treated by a routine second moment estimate, and the main task is to show that the expected value

$\displaystyle \sum_{n_1,\dots,n_k \in {\bf Z}^2} \mathop{\bf E}(\mu_{\delta,n_1} T \mu_{\delta,n_2} \dots T \mu_{\delta,n_k}) \ \ \ \ \ (3)$

becomes asymptotically independent of ${\delta}$ .

If all the ${n_1,\dots,n_k}$ were distinct then one could use independence to factor the expectation to get

$\displaystyle \sum_{n_1,\dots,n_k \in {\bf Z}^2} \mathop{\bf E}(\mu_{\delta,n_1}) T \mathop{\bf E}(\mu_{\delta,n_2}) \dots T \mathop{\bf E}(\mu_{\delta,n_k})$

which is a relatively straightforward expression to calculate (particularly in the model (1), where all the expectations here in fact vanish). The main difficulty is that there are a number of configurations in (3) in which various of the ${n_j}$ collide with each other, preventing one from easily factoring the expression. A typical problematic contribution for instance would be a sum of the form

$\displaystyle \sum_{n_1,n_2 \in {\bf Z}^2: n_1 \neq n_2} \mathop{\bf E}(\mu_{\delta,n_1} T \mu_{\delta,n_2} T \mu_{\delta,n_1} T \mu_{\delta,n_2}). \ \ \ \ \ (4)$

This is an example of what we call a non-split sum. This can be compared with the split sum

$\displaystyle \sum_{n_1,n_2 \in {\bf Z}^2: n_1 \neq n_2} \mathop{\bf E}(\mu_{\delta,n_1} T \mu_{\delta,n_1} T \mu_{\delta,n_2} T \mu_{\delta,n_2}). \ \ \ \ \ (5)$

If we ignore the constraint ${n_1 \neq n_2}$ in the latter sum, then it splits into

$\displaystyle f_\delta T g_\delta$

where

$\displaystyle f_\delta := \sum_{n_1 \in {\bf Z}^2} \mathop{\bf E}(\mu_{\delta,n_1} T \mu_{\delta,n_1})$

and

$\displaystyle g_\delta := \sum_{n_2 \in {\bf Z}^2} \mathop{\bf E}(\mu_{\delta,n_2} T \mu_{\delta,n_2})$

and one can hope to treat this sum by an induction hypothesis. (To actually deal with constraints such as ${n_1 \neq n_2}$ requires an inclusion-exclusion argument that creates some notational headaches but is ultimately manageable.) As the name suggests, the non-split configurations such as (4) cannot be factored in this fashion, and are the most difficult to handle. A direct computation using the triangle inequality (and a certain amount of combinatorics and induction) reveals that these sums are somewhat localised, in that dyadic portions such as

$\displaystyle \sum_{n_1,n_2 \in {\bf Z}^2: |n_1 - n_2| \sim R} \mathop{\bf E}(\mu_{\delta,n_1} T \mu_{\delta,n_2} T \mu_{\delta,n_1} T \mu_{\delta,n_2})$

exhibit power decay in ${R}$ (when measured in suitable function space norms), basically because of the large number of times one has to transition back and forth between ${n_1}$ and ${n_2}$ . Thus, morally at least, the dominant contribution to a non-split sum such as (4) comes from the local portion when ${n_2=n_1+O(1)}$ . From the translation and dilation invariance of ${T}$ this type of expression then simplifies to something like

$\displaystyle \varphi(z)^4 \sum_{n \in {\bf Z}^2} \eta( \frac{n-z}{\delta} )$

(plus negligible errors) for some reasonably decaying function ${\eta}$ , and this can be shown to converge to a weak limit as ${\delta \rightarrow 0}$ .

In principle all of these limits are computable, but the combinatorics is remarkably complicated, and while there is certainly some algebraic structure to the calculations, it does not seem to be easily describable in terms of an existing framework (e.g., that of free probability).

Local trigonometric polynomial approximations to the Riemann zeta function

3 May, 2019 in expository, math.CV, math.NT | Tags: polynomial approximation, Riemann zeta function, zero density theorems, zero-free region, zeroes | by Terence Tao | 29 comments

A useful rule of thumb in complex analysis is that holomorphic functions ${f(z)}$ behave like large degree polynomials ${P(z)}$ . This can be evidenced for instance at a “local” level by the Taylor series expansion for a complex analytic function in the disk, or at a “global” level by factorisation theorems such as the Weierstrass factorisation theorem (or the closely related Hadamard factorisation theorem). One can truncate these theorems in a variety of ways (e.g., Taylor’s theorem with remainder) to be able to approximate a holomorphic function by a polynomial on various domains.

In some cases it can be convenient instead to work with polynomials ${P(Z)}$ of another variable ${Z}$ such as ${Z = e^{2\pi i z}}$ (or more generally ${Z=e^{2\pi i z/N}}$ for a scaling parameter ${N}$ ). In the case of the Riemann zeta function, defined by meromorphic continuation of the formula

$\displaystyle \zeta(s) = \sum_{n=1}^\infty \frac{1}{n^s} \ \ \ \ \ (1)$

one ends up having the following heuristic approximation in the neighbourhood of a point ${\frac{1}{2}+it}$ on the critical line:

Heuristic 1 (Polynomial approximation) Let ${T \ggg 1}$ be a height, let ${t}$ be a “typical” element of ${[T,2T]}$ , and let ${1 \lll N \ll \log T}$ be an integer. Let ${\phi_t = \phi_{t,T}: {\bf C} \rightarrow {\bf C}}$ be the linear change of variables

$\displaystyle \phi_t(z) := \frac{1}{2} + it - \frac{2\pi i z}{\log T}.$

Then one has an approximation

$\displaystyle \zeta( \phi_t(z) ) \approx P_t( e^{2\pi i z/N} ) \ \ \ \ \ (2)$

for ${z = o(N)}$ and some polynomial ${P_t = P_{t,T}}$ of degree ${N}$ .

The requirement ${z=o(N)}$ is necessary since the right-hand side is periodic with period ${N}$ in the ${z}$ variable (or period ${\frac{2\pi i N}{\log T}}$ in the ${s = \phi_t(z)}$ variable), whereas the zeta function is not expected to have any such periodicity, even approximately.

Let us give two non-rigorous justifications of this heuristic. Firstly, it is standard that inside the critical strip (with ${\mathrm{Im}(s) = O(T)}$ ) we have an approximate form

$\displaystyle \zeta(s) \approx \sum_{n \leq T} \frac{1}{n^s}$

of (11). If we group the integers ${n}$ from ${1}$ to ${T}$ into ${N}$ bins depending on what powers of ${T^{1/N}}$ they lie between, we thus have

$\displaystyle \zeta(s) \approx \sum_{j=0}^N \sum_{T^{j/N} \leq n < T^{(j+1)/N}} \frac{1}{n^s}$

For ${s = \phi_t(z)}$ with ${z = o(N)}$ and ${T^{j/N} \leq n < T^{(j+1)/N}}$ we heuristically have

$\displaystyle \frac{1}{n^s} \approx \frac{1}{n^{\frac{1}{2}+it}} e^{2\pi i j z / N}$

and so

$\displaystyle \zeta(s) \approx \sum_{j=0}^N a_j(t) (e^{2\pi i z/N})^j$

where ${a_j(t)}$ are the partial Dirichlet series

$\displaystyle a_j(t) \approx \sum_{T^{j/N} \leq n < T^{(j+1)/N}} \frac{1}{n^{\frac{1}{2}+it}}. \ \ \ \ \ (3)$

This gives the desired polynomial approximation.

A second non-rigorous justification is as follows. From factorisation theorems such as the Hadamard factorisation theorem we expect to have

$\displaystyle \zeta(s) \propto \prod_\rho (s-\rho) \times \dots$

where ${\rho}$ runs over the non-trivial zeroes of ${\zeta}$ , and there are some additional factors arising from the trivial zeroes and poles of ${\zeta}$ which we will ignore here; we will also completely ignore the issue of how to renormalise the product to make it converge properly. In the region ${s = \frac{1}{2} + it + o( N / \log T) = \phi_t( \{ z: z = o(N) \})}$ , the dominant contribution to this product (besides multiplicative constants) should arise from zeroes ${\rho}$ that are also in this region. The Riemann-von Mangoldt formula suggests that for “typical” ${t}$ one should have about ${N}$ such zeroes. If one lets ${\rho_1,\dots,\rho_N}$ be any enumeration of ${N}$ zeroes closest to ${\frac{1}{2}+it}$ , and then repeats this set of zeroes periodically by period ${\frac{2\pi i N}{\log T}}$ , one then expects to have an approximation of the form

$\displaystyle \zeta(s) \propto \prod_{j=1}^N \prod_{k \in {\bf Z}} (s-(\rho_j+\frac{2\pi i kN}{\log T}) )$

again ignoring all issues of convergence. If one writes ${s = \phi_t(z)}$ and ${\rho_j = \phi_t(\lambda_j)}$ , then Euler’s famous product formula for sine basically gives

$\displaystyle \prod_{k \in {\bf Z}} (s-(\rho_j+\frac{2\pi i kN}{\log T}) ) \propto \prod_{k \in {\bf Z}} (z - (\lambda_j+2\pi k N) )$

$\displaystyle \propto (e^{2\pi i z/N} - e^{2\pi i \lambda j/N})$

(here we are glossing over some technical issues regarding renormalisation of the infinite products, which can be dealt with by studying the asymptotics as ${\mathrm{Im}(z) \rightarrow \infty}$ ) and hence we expect

$\displaystyle \zeta(s) \propto \prod_{j=1}^N (e^{2\pi i z/N} - e^{2\pi i \lambda j/N}).$

This again gives the desired polynomial approximation.

Below the fold we give a rigorous version of the second argument suitable for “microscale” analysis. More precisely, we will show

Theorem 2 Let ${N = N(T)}$ be an integer going sufficiently slowly to infinity. Let ${W_0 \ll N}$ go to zero sufficiently slowly depending on ${N}$ . Let ${t}$ be drawn uniformly at random from ${[T,2T]}$ . Then with probability ${1-o(1)}$ (in the limit ${T \rightarrow \infty}$ ), and possibly after adjusting ${N}$ by ${1}$ , there exists a polynomial ${P_t(Z)}$ of degree ${N}$ and obeying the functional equation (9) below, such that

$\displaystyle \zeta( \phi_t(z) ) = (1+o(1)) P_t( e^{2\pi i z/N} ) \ \ \ \ \ (4)$

whenever ${|z| \leq W_0}$ .

It should be possible to refine the arguments to extend this theorem to the mesoscale setting by letting ${N}$ be anything growing like ${o(\log T)}$ , and ${W_0}$ anything growing like ${o(N)}$ ; also we should be able to delete the need to adjust ${N}$ by ${1}$ . We have not attempted these optimisations here.

Many conjectures and arguments involving the Riemann zeta function can be heuristically translated into arguments involving the polynomials ${P_t(Z)}$ , which one can view as random degree ${N}$ polynomials if ${t}$ is interpreted as a random variable drawn uniformly at random from ${[T,2T]}$ . These can be viewed as providing a “toy model” for the theory of the Riemann zeta function, in which the complex analysis is simplified to the study of the zeroes and coefficients of this random polynomial (for instance, the role of the gamma function is now played by a monomial in ${Z}$ ). This model also makes the zeta function theory more closely resemble the function field analogues of this theory (in which the analogue of the zeta function is also a polynomial (or a rational function) in some variable ${Z}$ , as per the Weil conjectures). The parameter ${N}$ is at our disposal to choose, and reflects the scale ${\approx N/\log T}$ at which one wishes to study the zeta function. For “macroscopic” questions, at which one wishes to understand the zeta function at unit scales, it is natural to take ${N \approx \log T}$ (or very slightly larger), while for “microscopic” questions one would take ${N}$ close to ${1}$ and only growing very slowly with ${T}$ . For the intermediate “mesoscopic” scales one would take ${N}$ somewhere between ${1}$ and ${\log T}$ . Unfortunately, the statistical properties of ${P_t}$ are only understood well at a conjectural level at present; even if one assumes the Riemann hypothesis, our understanding of ${P_t}$ is largely restricted to the computation of low moments (e.g., the second or fourth moments) of various linear statistics of ${P_t}$ and related functions (e.g., ${1/P_t}$ , ${P'_t/P_t}$ , or ${\log P_t}$ ).

Let’s now heuristically explore the polynomial analogues of this theory in a bit more detail. The Riemann hypothesis basically corresponds to the assertion that all the ${N}$ zeroes of the polynomial ${P_t(Z)}$ lie on the unit circle ${|Z|=1}$ (which, after the change of variables ${Z = e^{2\pi i z/N}}$ , corresponds to ${z}$ being real); in a similar vein, the GUE hypothesis corresponds to ${P_t(Z)}$ having the asymptotic law of a random scalar ${a_N(t)}$ times the characteristic polynomial of a random unitary ${N \times N}$ matrix. Next, we consider what happens to the functional equation

$\displaystyle \zeta(s) = \chi(s) \zeta(1-s) \ \ \ \ \ (5)$

where

$\displaystyle \chi(s) := 2^s \pi^{s-1} \sin(\frac{\pi s}{2}) \Gamma(1-s).$

A routine calculation involving Stirling’s formula reveals that

$\displaystyle \chi(\frac{1}{2}+it) = (1+o(1)) e^{-2\pi i L(t)} \ \ \ \ \ (6)$

with ${L(t) := \frac{t}{2\pi} \log \frac{t}{2\pi} - \frac{t}{2\pi} + \frac{7}{8}}$ ; one also has the closely related approximation

$\displaystyle \frac{\chi'}{\chi}(s) = -\log T + O(1) \ \ \ \ \ (7)$

and hence

$\displaystyle \chi(\phi_t(z)) = (1+o(1)) e^{-2\pi i \theta(t)} e^{2\pi i z} \ \ \ \ \ (8)$

when ${z = o(\log T)}$ . Since ${\zeta(1-s) = \overline{\zeta(\overline{1-s})}}$ , applying (5) with ${s = \phi_t(z)}$ and using the approximation (2) suggests a functional equation for ${P_t}$ :

$\displaystyle P_t(e^{2\pi i z/N}) = e^{-2\pi i L(t)} e^{2\pi i z} \overline{P_t(e^{2\pi i \overline{z}/N})}$

or in terms of ${Z := e^{2\pi i z/N}}$ ,

$\displaystyle P_t(Z) = e^{-2\pi i L(t)} Z^N \overline{P_t}(1/Z) \ \ \ \ \ (9)$

where ${\overline{P_t}(Z) := \overline{P_t(\overline{Z})}}$ is the polynomial ${P_t}$ with all the coefficients replaced by their complex conjugate. Thus if we write

$\displaystyle P_t(Z) = \sum_{j=0}^N a_j Z^j$

then the functional equation can be written as

$\displaystyle a_j(t) = e^{-2\pi i L(t)} \overline{a_{N-j}(t)}.$

We remark that if we use the heuristic (3) (interpreting the cutoffs in the ${n}$ summation in a suitably vague fashion) then this equation can be viewed as an instance of the Poisson summation formula.

Another consequence of the functional equation is that the zeroes of ${P_t}$ are symmetric with respect to inversion ${Z \mapsto 1/\overline{Z}}$ across the unit circle. This is of course consistent with the Riemann hypothesis, but does not obviously imply it. The phase ${L(t)}$ is of little consequence in this functional equation; one could easily conceal it by working with the phase rotation ${e^{\pi i L(t)} P_t}$ of ${P_t}$ instead.

One consequence of the functional equation is that ${e^{\pi i L(t)} e^{-i N \theta/2} P_t(e^{i\theta})}$ is real for any ${\theta \in {\bf R}}$ ; the same is then true for the derivative ${e^{\pi i L(t)} e^{i N \theta} (i e^{i\theta} P'_t(e^{i\theta}) - i \frac{N}{2} P_t(e^{i\theta})}$ . Among other things, this implies that ${P'_t(e^{i\theta})}$ cannot vanish unless ${P_t(e^{i\theta})}$ does also; thus the zeroes of ${P'_t}$ will not lie on the unit circle except where ${P_t}$ has repeated zeroes. The analogous statement is true for ${\zeta}$ ; the zeroes of ${\zeta'}$ will not lie on the critical line except where ${\zeta}$ has repeated zeroes.

Relating to this fact, it is a classical result of Speiser that the Riemann hypothesis is true if and only if all the zeroes of the derivative ${\zeta'}$ of the zeta function in the critical strip lie on or to the right of the critical line. The analogous result for polynomials is

Proposition 3 We have

$\displaystyle \# \{ |Z| = 1: P_t(Z) = 0 \} = N - 2 \# \{ |Z| > 1: P'_t(Z) = 0 \}$

(where all zeroes are counted with multiplicity.) In particular, the zeroes of ${P_t(Z)}$ all lie on the unit circle if and only if the zeroes of ${P'_t(Z)}$ lie in the closed unit disk.

Proof: From the functional equation we have

$\displaystyle \# \{ |Z| = 1: P_t(Z) = 0 \} = N - 2 \# \{ |Z| > 1: P_t(Z) = 0 \}.$

Thus it will suffice to show that ${P_t}$ and ${P'_t}$ have the same number of zeroes outside the closed unit disk.

Set ${f(z) := z \frac{P'(z)}{P(z)}}$ , then ${f}$ is a rational function that does not have a zero or pole at infinity. For ${e^{i\theta}}$ not a zero of ${P_t}$ , we have already seen that ${e^{\pi i L(t)} e^{-i N \theta/2} P_t(e^{i\theta})}$ and ${e^{\pi i L(t)} e^{i N \theta} (i e^{i\theta} P'_t(e^{i\theta}) - i \frac{N}{2} P_t(e^{i\theta})}$ are real, so on dividing we see that ${i f(e^{i\theta}) - \frac{iN}{2}}$ is always real, that is to say

$\displaystyle \mathrm{Re} f(e^{i\theta}) = \frac{N}{2}.$

(This can also be seen by writing ${f(e^{i\theta}) = \sum_\lambda \frac{1}{1-e^{-i\theta} \lambda}}$ , where ${\lambda}$ runs over the zeroes of ${P_t}$ , and using the fact that these zeroes are symmetric with respect to reflection across the unit circle.) When ${e^{i\theta}}$ is a zero of ${P_t}$ , ${f(z)}$ has a simple pole at ${e^{i\theta}}$ with residue a positive multiple of ${e^{i\theta}}$ , and so ${f(z)}$ stays on the right half-plane if one traverses a semicircular arc around ${e^{i\theta}}$ outside the unit disk. From this and continuity we see that ${f}$ stays on the right-half plane in a circle slightly larger than the unit circle, and hence by the argument principle it has the same number of zeroes and poles outside of this circle, giving the claim. $\Box$

From the functional equation and the chain rule, ${Z}$ is a zero of ${P'_t}$ if and only if ${1/\overline{Z}}$ is a zero of ${N P_t - P'_t}$ . We can thus write the above proposition in the equivalent form

$\displaystyle \# \{ |Z| = 1: P_t(Z) = 0 \} = N - 2 \# \{ |Z| < 1: NP_t(Z) - P'_t(Z) = 0 \}.$

One can use this identity to get a lower bound on the number of zeroes of ${P_t}$ by the method of mollifiers. Namely, for any other polynomial ${M_t}$ , we clearly have

$\displaystyle \# \{ |Z| = 1: P_t(Z) = 0 \}$

$\displaystyle \geq N - 2 \# \{ |Z| < 1: M_t(Z)(NP_t(Z) - P'_t(Z)) = 0 \}.$

By Jensen’s formula, we have for any ${r>1}$ that

$\displaystyle \log |M_t(0)| |NP_t(0)-P'_t(0)|$

$\displaystyle \leq -(\log r) \# \{ |Z| < 1: M_t(Z)(NP_t(Z) - P'_t(Z)) = 0 \}$

$\displaystyle + \frac{1}{2\pi} \int_0^{2\pi} \log |M_t(re^{i\theta})(NP_t(e^{i\theta}) - P'_t(re^{i\theta}))|\ d\theta.$

We therefore have

$\displaystyle \# \{ |Z| = 1: P_t(Z) = 0 \} \geq N + \frac{2}{\log r} \log |M_t(0)| |NP_t(0)-P'_t(0)|$

$\displaystyle - \frac{1}{\log r} \frac{1}{2\pi} \int_0^{2\pi} \log |M_t(re^{i\theta})(NP_t(e^{i\theta}) - P'_t(re^{i\theta}))|^2\ d\theta.$

As the logarithm function is concave, we can apply Jensen’s inequality to conclude

$\displaystyle {\bf E} \# \{ |Z| = 1: P_t(Z) = 0 \} \geq N$

$\displaystyle + {\bf E} \frac{2}{\log r} \log |M_t(0)| |NP_t(0)-P'_t(0)|$

$\displaystyle - \frac{1}{\log r} \log \left( \frac{1}{2\pi} \int_0^{2\pi} {\bf E} |M_t(re^{i\theta})(NP_t(e^{i\theta}) - P'_t(re^{i\theta}))|^2\ d\theta\right).$

where the expectation is over the ${t}$ parameter. It turns out that by choosing the mollifier ${M_t}$ carefully in order to make ${M_t P_t}$ behave like the function ${1}$ (while keeping the degree ${M_t}$ small enough that one can compute the second moment here), and then optimising in ${r}$ , one can use this inequality to get a positive fraction of zeroes of ${P_t}$ on the unit circle on average. This is the polynomial analogue of a classical argument of Levinson, who used this to show that at least one third of the zeroes of the Riemann zeta function are on the critical line; all later improvements on this fraction have been based on some version of Levinson’s method, mainly focusing on more advanced choices for the mollifier ${M_t}$ and of the differential operator ${N - \partial_z}$ that implicitly appears in the above approach. (The most recent lower bound I know of is ${0.4191637}$ , due to Pratt and Robles. In principle (as observed by Farmer) this bound can get arbitrarily close to ${1}$ if one is allowed to use arbitrarily long mollifiers, but establishing this seems of comparable difficulty to unsolved problems such as the pair correlation conjecture; see this paper of Radziwill for more discussion.) A variant of these techniques can also establish “zero density estimates” of the following form: for any ${W \geq 1}$ , the number of zeroes of ${P_t}$ that lie further than ${\frac{W}{N}}$ from the unit circle is of order ${O( e^{-cW} N )}$ on average for some absolute constant ${c>0}$ . Thus, roughly speaking, most zeroes of ${P_t}$ lie within ${O(1/N)}$ of the unit circle. (Analogues of these results for the Riemann zeta function were worked out by Selberg, by Jutila, and by Conrey, with increasingly strong values of ${c}$ .)

The zeroes of ${P'_t}$ tend to live somewhat closer to the origin than the zeroes of ${P_t}$ . Suppose for instance that we write

$\displaystyle P_t(Z) = \sum_{j=0}^N a_j(t) Z^j = a_N(t) \prod_{j=1}^N (Z - \lambda_j)$

where ${\lambda_1,\dots,\lambda_N}$ are the zeroes of ${P_t(Z)}$ , then by evaluating at zero we see that

$\displaystyle \lambda_1 \dots \lambda_N = (-1)^N a_0(t) / a_N(t)$

and the right-hand side is of unit magnitude by the functional equation. However, if we differentiate

$\displaystyle P'_t(Z) = \sum_{j=1}^N a_j(t) j Z^{j-1} = N a_N(t) \prod_{j=1}^{N-1} (Z - \lambda'_j)$

where ${\lambda'_1,\dots,\lambda'_{N-1}}$ are the zeroes of ${P'_t}$ , then by evaluating at zero we now see that

$\displaystyle \lambda'_1 \dots \lambda'_{N-1} = (-1)^N a_1(t) / N a_N(t).$

The right-hand side would now be typically expected to be of size ${O(1/N) \approx \exp(- \log N)}$ , and so on average we expect the ${\lambda'_j}$ to have magnitude like ${\exp( - \frac{\log N}{N} )}$ , that is to say pushed inwards from the unit circle by a distance roughly ${\frac{\log N}{N}}$ . The analogous result for the Riemann zeta function is that the zeroes of ${\zeta'(s)}$ at height ${\sim T}$ lie at a distance roughly ${\frac{\log\log T}{\log T}}$ to the right of the critical line on the average; see this paper of Levinson and Montgomery for a precise statement.

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246C notes 4: Brownian motion, conformal invariance, and SLE

29 May, 2018 in 246C - complex analysis, math.CV, math.PR | Tags: Brownian motion, conformal mapping, Schramm-Loewner evolution | by Terence Tao | 6 comments

Important note: As this is not a course in probability, we will try to avoid developing the general theory of stochastic calculus (which includes such concepts as filtrations, martingales, and Ito calculus). This will unfortunately limit what we can actually prove rigorously, and so at some places the arguments will be somewhat informal in nature. A rigorous treatment of many of the topics here can be found for instance in Lawler’s Conformally Invariant Processes in the Plane, from which much of the material here is drawn.

In these notes, random variables will be denoted in boldface.

Definition 1 A real random variable ${\mathbf{X}}$ is said to be normally distributed with mean ${x_0 \in {\bf R}}$ and variance ${\sigma^2 > 0}$ if one has

$\displaystyle \mathop{\bf E} F(\mathbf{X}) = \frac{1}{\sqrt{2\pi} \sigma} \int_{\bf R} e^{-(x-x_0)^2/2\sigma^2} F(x)\ dx$

for all test functions ${F \in C_c({\bf R})}$ . Similarly, a complex random variable ${\mathbf{Z}}$ is said to be normally distributed with mean ${z_0 \in {\bf R}}$ and variance ${\sigma^2>0}$ if one has

$\displaystyle \mathop{\bf E} F(\mathbf{Z}) = \frac{1}{\pi \sigma^2} \int_{\bf C} e^{-|z-x_0|^2/\sigma^2} F(z)\ dx dy$

for all test functions ${F \in C_c({\bf C})}$ , where ${dx dy}$ is the area element on ${{\bf C}}$ .

A real Brownian motion with base point ${x_0 \in {\bf R}}$ is a random, almost surely continuous function ${\mathbf{B}^{x_0}: [0,+\infty) \rightarrow {\bf R}}$ (using the locally uniform topology on continuous functions) with the property that (almost surely) ${\mathbf{B}^{x_0}(0) = x_0}$ , and for any sequence of times ${0 \leq t_0 < t_1 < t_2 < \dots < t_n}$ , the increments ${\mathbf{B}^{x_0}(t_i) - \mathbf{B}^{x_0}(t_{i-1})}$ for ${i=1,\dots,n}$ are independent real random variables that are normally distributed with mean zero and variance ${t_i - t_{i-1}}$ . Similarly, a complex Brownian motion with base point ${z_0 \in {\bf R}}$ is a random, almost surely continuous function ${\mathbf{B}^{z_0}: [0,+\infty) \rightarrow {\bf R}}$ with the property that ${\mathbf{B}^{z_0}(0) = z_0}$ and for any sequence of times ${0 \leq t_0 < t_1 < t_2 < \dots < t_n}$ , the increments ${\mathbf{B}^{z_0}(t_i) - \mathbf{B}^{z_0}(t_{i-1})}$ for ${i=1,\dots,n}$ are independent complex random variables that are normally distributed with mean zero and variance ${t_i - t_{i-1}}$ .

Remark 2 Thanks to the central limit theorem, the hypothesis that the increments ${\mathbf{B}^{x_0}(t_i) - \mathbf{B}^{x_0}(t_{i-1})}$ be normally distributed can be dropped from the definition of a Brownian motion, so long as one retains the independence and the normalisation of the mean and variance (technically one also needs some uniform integrability on the increments beyond the second moment, but we will not detail this here). A similar statement is also true for the complex Brownian motion (where now we need to normalise the variances and covariances of the real and imaginary parts of the increments).

Real and complex Brownian motions exist from any base point ${x_0}$ or ${z_0}$ ; see e.g. this previous blog post for a construction. We have the following simple invariances:

Exercise 3

(i) (Translation invariance) If ${\mathbf{B}^{x_0}}$ is a real Brownian motion with base point ${x_0 \in {\bf R}}$ , and ${h \in {\bf R}}$ , show that ${\mathbf{B}^{x_0}+h}$ is a real Brownian motion with base point ${x_0+h}$ . Similarly, if ${\mathbf{B}^{z_0}}$ is a complex Brownian motion with base point ${z_0 \in {\bf R}}$ , and ${h \in {\bf C}}$ , show that ${\mathbf{B}^{z_0}+c}$ is a complex Brownian motion with base point ${z_0+h}$ .

(ii) (Dilation invariance) If ${\mathbf{B}^{0}}$ is a real Brownian motion with base point ${0}$ , and ${\lambda \in {\bf R}}$ is non-zero, show that ${t \mapsto \lambda \mathbf{B}^0(t / |\lambda|^{1/2})}$ is also a real Brownian motion with base point ${0}$ . Similarly, if ${\mathbf{B}^0}$ is a complex Brownian motion with base point ${0}$ , and ${\lambda \in {\bf C}}$ is non-zero, show that ${t \mapsto \lambda \mathbf{B}^0(t / |\lambda|^{1/2})}$ is also a complex Brownian motion with base point ${0}$ .

(iii) (Real and imaginary parts) If ${\mathbf{B}^0}$ is a complex Brownian motion with base point ${0}$ , show that ${\sqrt{2} \mathrm{Re} \mathbf{B}^0}$ and ${\sqrt{2} \mathrm{Im} \mathbf{B}^0}$ are independent real Brownian motions with base point ${0}$ . Conversely, if ${\mathbf{B}^0_1, \mathbf{B}^0_2}$ are independent real Brownian motions of base point ${0}$ , show that ${\frac{1}{\sqrt{2}} (\mathbf{B}^0_1 + i \mathbf{B}^0_2)}$ is a complex Brownian motion with base point ${0}$ .

The next lemma is a special case of the optional stopping theorem.

Lemma 4 (Optional stopping identities)

(i) (Real case) Let ${\mathbf{B}^{x_0}}$ be a real Brownian motion with base point ${x_0 \in {\bf R}}$ . Let ${\mathbf{t}}$ be a bounded stopping time – a bounded random variable with the property that for any time ${t \geq 0}$ , the event that ${\mathbf{t} \leq t}$ is determined by the values of the trajectory ${\mathbf{B}^{x_0}}$ for times up to ${t}$ (or more precisely, this event is measurable with respect to the ${\sigma}$ algebra generated by this proprtion of the trajectory). Then
$\displaystyle \mathop{\bf E} \mathbf{B}^{x_0}(\mathbf{t}) = x_0$

and

$\displaystyle \mathop{\bf E} (\mathbf{B}^{x_0}(\mathbf{t})-x_0)^2 - \mathbf{t} = 0$

and

$\displaystyle \mathop{\bf E} (\mathbf{B}^{x_0}(\mathbf{t})-x_0)^4 = O( \mathop{\bf E} \mathbf{t}^2 ).$

(ii) (Complex case) Let ${\mathbf{B}^{z_0}}$ be a real Brownian motion with base point ${z_0 \in {\bf R}}$ . Let ${\mathbf{t}}$ be a bounded stopping time – a bounded random variable with the property that for any time ${t \geq 0}$ , the event that ${\mathbf{t} \leq t}$ is determined by the values of the trajectory ${\mathbf{B}^{x_0}}$ for times up to ${t}$ . Then
$\displaystyle \mathop{\bf E} \mathbf{B}^{z_0}(\mathbf{t}) = z_0$

$\displaystyle \mathop{\bf E} (\mathrm{Re}(\mathbf{B}^{z_0}(\mathbf{t})-z_0))^2 - \frac{1}{2} \mathbf{t} = 0$

$\displaystyle \mathop{\bf E} (\mathrm{Im}(\mathbf{B}^{z_0}(\mathbf{t})-z_0))^2 - \frac{1}{2} \mathbf{t} = 0$

$\displaystyle \mathop{\bf E} \mathrm{Re}(\mathbf{B}^{z_0}(\mathbf{t})-z_0) \mathrm{Im}(\mathbf{B}^{z_0}(\mathbf{t})-z_0) = 0$

$\displaystyle \mathop{\bf E} |\mathbf{B}^{x_0}(\mathbf{t})-z_0|^4 = O( \mathop{\bf E} \mathbf{t}^2 ).$

Proof: (Slightly informal) We just prove (i) and leave (ii) as an exercise. By translation invariance we can take ${x_0=0}$ . Let ${T}$ be an upper bound for ${\mathbf{t}}$ . Since ${\mathbf{B}^0(T)}$ is a real normally distributed variable with mean zero and variance ${T}$ , we have

$\displaystyle \mathop{\bf E} \mathbf{B}^0( T ) = 0$

and

$\displaystyle \mathop{\bf E} \mathbf{B}^0( T )^2 = T$

and

$\displaystyle \mathop{\bf E} \mathbf{B}^0( T )^4 = 3T^2.$

By the law of total expectation, we thus have

$\displaystyle \mathop{\bf E} \mathop{\bf E}(\mathbf{B}^0( T ) | \mathbf{t}, \mathbf{B}^{z_0}(\mathbf{t}) ) = 0$

and

$\displaystyle \mathop{\bf E} \mathop{\bf E}((\mathbf{B}^0( T ))^2 | \mathbf{t}, \mathbf{B}^{z_0}(\mathbf{t}) ) = T$

and

$\displaystyle \mathop{\bf E} \mathop{\bf E}((\mathbf{B}^0( T ))^4 | \mathbf{t}, \mathbf{B}^{z_0}(\mathbf{t}) ) = 3T^2$

where the inner conditional expectations are with respect to the event that ${\mathbf{t}, \mathbf{B}^{0}(\mathbf{t})}$ attains a particular point in ${S}$ . However, from the independent increment nature of Brownian motion, once one conditions ${(\mathbf{t}, \mathbf{B}^{0}(\mathbf{t}))}$ to a fixed point ${(t, x)}$ , the random variable ${\mathbf{B}^0(T)}$ becomes a real normally distributed variable with mean ${x}$ and variance ${T-t}$ . Thus we have

$\displaystyle \mathop{\bf E}(\mathbf{B}^0( T ) | \mathbf{t}, \mathbf{B}^{z_0}(\mathbf{t}) ) = \mathbf{B}^{z_0}(\mathbf{t})$

and

$\displaystyle \mathop{\bf E}( (\mathbf{B}^0( T ))^2 | \mathbf{t}, \mathbf{B}^{z_0}(\mathbf{t}) ) = \mathbf{B}^{z_0}(\mathbf{t})^2 + T - \mathbf{t}$

and

$\displaystyle \mathop{\bf E}( (\mathbf{B}^0( T ))^4 | \mathbf{t}, \mathbf{B}^{z_0}(\mathbf{t}) ) = \mathbf{B}^{z_0}(\mathbf{t})^4 + 6(T - \mathbf{t}) \mathbf{B}^{z_0}(\mathbf{t})^2 + 3(T - \mathbf{t})^2$

which give the first two claims, and (after some algebra) the identity

$\displaystyle \mathop{\bf E} \mathbf{B}^{z_0}(\mathbf{t})^4 - 6 \mathbf{t} \mathbf{B}^{z_0}(\mathbf{t})^2 + 3 \mathbf{t}^2 = 0$

which then also gives the third claim. $\Box$

Exercise 5 Prove the second part of Lemma 4.

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246C notes 3: Univalent functions, the Loewner equation, and the Bieberbach conjecture

2 May, 2018 in 246C - complex analysis, math.AP, math.CA, math.CV | Tags: Bieberbach conjecture, Loewner equation, Univalent functions | by Terence Tao | 13 comments

We now approach conformal maps from yet another perspective. Given an open subset ${U}$ of the complex numbers ${{\bf C}}$ , define a univalent function on ${U}$ to be a holomorphic function ${f: U \rightarrow {\bf C}}$ that is also injective. We will primarily be studying this concept in the case when ${U}$ is the unit disk ${D(0,1) := \{ z \in {\bf C}: |z| < 1 \}}$ .

Clearly, a univalent function ${f: D(0,1) \rightarrow {\bf C}}$ on the unit disk is a conformal map from ${D(0,1)}$ to the image ${f(D(0,1))}$ ; in particular, ${f(D(0,1))}$ is simply connected, and not all of ${{\bf C}}$ (since otherwise the inverse map ${f^{-1}: {\bf C} \rightarrow D(0,1)}$ would violate Liouville’s theorem). In the converse direction, the Riemann mapping theorem tells us that every open simply connected proper subset ${V \subsetneq {\bf C}}$ of the complex numbers is the image of a univalent function on ${D(0,1)}$ . Furthermore, if ${V}$ contains the origin, then the univalent function ${f: D(0,1) \rightarrow {\bf C}}$ with this image becomes unique once we normalise ${f(0) = 0}$ and ${f'(0) > 0}$ . Thus the Riemann mapping theorem provides a one-to-one correspondence between open simply connected proper subsets of the complex plane containing the origin, and univalent functions ${f: D(0,1) \rightarrow {\bf C}}$ with ${f(0)=0}$ and ${f'(0)>0}$ . We will focus particular attention on the univalent functions ${f: D(0,1) \rightarrow {\bf C}}$ with the normalisation ${f(0)=0}$ and ${f'(0)=1}$ ; such functions will be called schlicht functions.

One basic example of a univalent function on ${D(0,1)}$ is the Cayley transform ${z \mapsto \frac{1+z}{1-z}}$ , which is a Möbius transformation from ${D(0,1)}$ to the right half-plane ${\{ \mathrm{Re}(z) > 0 \}}$ . (The slight variant ${z \mapsto \frac{1-z}{1+z}}$ is also referred to as the Cayley transform, as is the closely related map ${z \mapsto \frac{z-i}{z+i}}$ , which maps ${D(0,1)}$ to the upper half-plane.) One can square this map to obtain a further univalent function ${z \mapsto \left( \frac{1+z}{1-z} \right)^2}$ , which now maps ${D(0,1)}$ to the complex numbers with the negative real axis ${(-\infty,0]}$ removed. One can normalise this function to be schlicht to obtain the Koebe function

$\displaystyle f(z) := \frac{1}{4}\left( \left( \frac{1+z}{1-z} \right)^2 - 1\right) = \frac{z}{(1-z)^2}, \ \ \ \ \ (1)$

which now maps ${D(0,1)}$ to the complex numbers with the half-line ${(-\infty,-1/4]}$ removed. A little more generally, for any ${\theta \in {\bf R}}$ we have the rotated Koebe function

$\displaystyle f(z) := \frac{z}{(1 - e^{i\theta} z)^2} \ \ \ \ \ (2)$

that is a schlicht function that maps ${D(0,1)}$ to the complex numbers with the half-line ${\{ -re^{-i\theta}: r \geq 1/4\}}$ removed.

Every schlicht function ${f: D(0,1) \rightarrow {\bf C}}$ has a convergent Taylor expansion

$\displaystyle f(z) = a_1 z + a_2 z^2 + a_3 z^3 + \dots$

for some complex coefficients ${a_1,a_2,\dots}$ with ${a_1=1}$ . For instance, the Koebe function has the expansion

$\displaystyle f(z) = z + 2 z^2 + 3 z^3 + \dots = \sum_{n=1}^\infty n z^n$

and similarly the rotated Koebe function has the expansion

$\displaystyle f(z) = z + 2 e^{i\theta} z^2 + 3 e^{2i\theta} z^3 + \dots = \sum_{n=1}^\infty n e^{(n-1)\theta} z^n.$

Intuitively, the Koebe function and its rotations should be the “largest” schlicht functions available. This is formalised by the famous Bieberbach conjecture, which asserts that for any schlicht function, the coefficients ${a_n}$ should obey the bound ${|a_n| \leq n}$ for all ${n}$ . After a large number of partial results, this conjecture was eventually solved by de Branges; see for instance this survey of Korevaar or this survey of Koepf for a history.

It turns out that to resolve these sorts of questions, it is convenient to restrict attention to schlicht functions ${g: D(0,1) \rightarrow {\bf C}}$ that are odd, thus ${g(-z)=-g(z)}$ for all ${z}$ , and the Taylor expansion now reads

$\displaystyle g(z) = b_1 z + b_3 z^3 + b_5 z^5 + \dots$

for some complex coefficients ${b_1,b_3,\dots}$ with ${b_1=1}$ . One can transform a general schlicht function ${f: D(0,1) \rightarrow {\bf C}}$ to an odd schlicht function ${g: D(0,1) \rightarrow {\bf C}}$ by observing that the function ${f(z^2)/z^2: D(0,1) \rightarrow {\bf C}}$ , after removing the singularity at zero, is a non-zero function that equals ${1}$ at the origin, and thus (as ${D(0,1)}$ is simply connected) has a unique holomorphic square root ${(f(z^2)/z^2)^{1/2}}$ that also equals ${1}$ at the origin. If one then sets

$\displaystyle g(z) := z (f(z^2)/z^2)^{1/2} \ \ \ \ \ (3)$

it is not difficult to verify that ${g}$ is an odd schlicht function which additionally obeys the equation

$\displaystyle f(z^2) = g(z)^2. \ \ \ \ \ (4)$

Conversely, given an odd schlicht function ${g}$ , the formula (4) uniquely determines a schlicht function ${f}$ .

For instance, if ${f}$ is the Koebe function (1), ${g}$ becomes

$\displaystyle g(z) = \frac{z}{1-z^2} = z + z^3 + z^5 + \dots, \ \ \ \ \ (5)$

which maps ${D(0,1)}$ to the complex numbers with two slits ${\{ \pm iy: y > 1/2 \}}$ removed, and if ${f}$ is the rotated Koebe function (2), ${g}$ becomes

$\displaystyle g(z) = \frac{z}{1- e^{i\theta} z^2} = z + e^{i\theta} z^3 + e^{2i\theta} z^5 + \dots. \ \ \ \ \ (6)$

De Branges established the Bieberbach conjecture by first proving an analogous conjecture for odd schlicht functions known as Robertson’s conjecture. More precisely, we have

Theorem 1 (de Branges’ theorem) Let ${n \geq 1}$ be a natural number.

(i) (Robertson conjecture) If ${g(z) = b_1 z + b_3 z^3 + b_5 z^5 + \dots}$ is an odd schlicht function, then
$\displaystyle \sum_{k=1}^n |b_{2k-1}|^2 \leq n.$

(ii) (Bieberbach conjecture) If ${f(z) = a_1 z + a_2 z^2 + a_3 z^3 + \dots}$ is a schlicht function, then
$\displaystyle |a_n| \leq n.$

It is easy to see that the Robertson conjecture for a given value of ${n}$ implies the Bieberbach conjecture for the same value of ${n}$ . Indeed, if ${f(z) = a_1 z + a_2 z^2 + a_3 z^3 + \dots}$ is schlicht, and ${g(z) = b_1 z + b_3 z^3 + b_5 z^5 + \dots}$ is the odd schlicht function given by (3), then from extracting the ${z^{2n}}$ coefficient of (4) we obtain a formula

$\displaystyle a_n = \sum_{j=1}^n b_{2j-1} b_{2(n+1-j)-1}$

for the coefficients of ${f}$ in terms of the coefficients of ${g}$ . Applying the Cauchy-Schwarz inequality, we derive the Bieberbach conjecture for this value of ${n}$ from the Robertson conjecture for the same value of ${n}$ . We remark that Littlewood and Paley had conjectured a stronger form ${|b_{2k-1}| \leq 1}$ of Robertson’s conjecture, but this was disproved for ${k=3}$ by Fekete and Szegö.

To prove the Robertson and Bieberbach conjectures, one first takes a logarithm and deduces both conjectures from a similar conjecture about the Taylor coefficients of ${\log \frac{f(z)}{z}}$ , known as the Milin conjecture. Next, one continuously enlarges the image ${f(D(0,1))}$ of the schlicht function to cover all of ${{\bf C}}$ ; done properly, this places the schlicht function ${f}$ as the initial function ${f = f_0}$ in a sequence ${(f_t)_{t \geq 0}}$ of univalent maps ${f_t: D(0,1) \rightarrow {\bf C}}$ known as a Loewner chain. The functions ${f_t}$ obey a useful differential equation known as the Loewner equation, that involves an unspecified forcing term ${\mu_t}$ (or ${\theta(t)}$ , in the case that the image is a slit domain) coming from the boundary; this in turn gives useful differential equations for the Taylor coefficients of ${f(z)}$ , ${g(z)}$ , or ${\log \frac{f(z)}{z}}$ . After some elementary calculus manipulations to “integrate” this equations, the Bieberbach, Robertson, and Milin conjectures are then reduced to establishing the non-negativity of a certain explicit hypergeometric function, which is non-trivial to prove (and will not be done here, except for small values of ${n}$ ) but for which several proofs exist in the literature.

The theory of Loewner chains subsequently became fundamental to a more recent topic in complex analysis, that of the Schramm-Loewner equation (SLE), which is the focus of the next and final set of notes.

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246C notes 2: Circle packings, conformal maps, and quasiconformal maps

12 April, 2018 in 246C - complex analysis, math.CV, math.MG | Tags: circle packing, quasiconformal mapping, Riemann mapping theorem | by Terence Tao | 24 comments

We now leave the topic of Riemann surfaces, and turn now to the (loosely related) topic of conformal mapping (and quasiconformal mapping). Recall that a conformal map ${f: U \rightarrow V}$ from an open subset ${U}$ of the complex plane to another open set ${V}$ is a map that is holomorphic and bijective, which (by Rouché’s theorem) also forces the derivative of ${f}$ to be nowhere vanishing. We then say that the two open sets ${U,V}$ are conformally equivalent. From the Cauchy-Riemann equations we see that conformal maps are orientation-preserving and angle-preserving; from the Newton approximation ${f( z_0 + \Delta z) \approx f(z_0) + f'(z_0) \Delta z + O( |\Delta z|^2)}$ we see that they almost preserve small circles, indeed for ${\varepsilon}$ small the circle ${\{ z: |z-z_0| = \varepsilon\}}$ will approximately map to ${\{ w: |w - f(z_0)| = |f'(z_0)| \varepsilon \}}$ .

In previous quarters, we proved a fundamental theorem about this concept, the Riemann mapping theorem:

Theorem 1 (Riemann mapping theorem) Let ${U}$ be a simply connected open subset of ${{\bf C}}$ that is not all of ${{\bf C}}$ . Then ${U}$ is conformally equivalent to the unit disk ${D(0,1)}$ .

This theorem was proven in these 246A lecture notes, using an argument of Koebe. At a very high level, one can sketch Koebe’s proof of the Riemann mapping theorem as follows: among all the injective holomorphic maps ${f: U \rightarrow D(0,1)}$ from ${U}$ to ${D(0,1)}$ that map some fixed point ${z_0 \in U}$ to ${0}$ , pick one that maximises the magnitude ${|f'(z_0)|}$ of the derivative (ignoring for this discussion the issue of proving that a maximiser exists). If ${f(U)}$ avoids some point in ${D(0,1)}$ , one can compose ${f}$ with various holomorphic maps and use Schwarz’s lemma and the chain rule to increase ${|f'(z_0)|}$ without destroying injectivity; see the previous lecture notes for details. The conformal map ${\phi: U \rightarrow D(0,1)}$ is unique up to Möbius automorphisms of the disk; one can fix the map by picking two distinct points ${z_0,z_1}$ in ${U}$ , and requiring ${\phi(z_0)}$ to be zero and ${\phi(z_1)}$ to be positive real.

It is a beautiful observation of Thurston that the concept of a conformal mapping has a discrete counterpart, namely the mapping of one circle packing to another. Furthermore, one can run a version of Koebe’s argument (using now a discrete version of Perron’s method) to prove the Riemann mapping theorem through circle packings. In principle, this leads to a mostly elementary approach to conformal geometry, based on extremely classical mathematics that goes all the way back to Apollonius. However, in order to prove the basic existence and uniqueness theorems of circle packing, as well as the convergence to conformal maps in the continuous limit, it seems to be necessary (or at least highly convenient) to use much more modern machinery, including the theory of quasiconformal mapping, and also the Riemann mapping theorem itself (so in particular we are not structuring these notes to provide a completely independent proof of that theorem, though this may well be possible).

To make the above discussion more precise we need some notation.

Definition 2 (Circle packing) A (finite) circle packing is a finite collection ${(C_j)_{j \in J}}$ of circles ${C_j = \{ z \in {\bf C}: |z-z_j| = r_j\}}$ in the complex numbers indexed by some finite set ${J}$ , whose interiors are all disjoint (but which are allowed to be tangent to each other), and whose union is connected. The nerve of a circle packing is the finite graph whose vertices ${\{z_j: j \in J \}}$ are the centres of the circle packing, with two such centres connected by an edge if the circles are tangent. (In these notes all graphs are undirected, finite and simple, unless otherwise specified.)

It is clear that the nerve of a circle packing is connected and planar, since one can draw the nerve by placing each vertex (tautologically) in its location in the complex plane, and drawing each edge by the line segment between the centres of the circles it connects (this line segment will pass through the point of tangency of the two circles). Later in these notes we will also have to consider some infinite circle packings, most notably the infinite regular hexagonal circle packing.

The first basic theorem in the subject is the following converse statement:

Theorem 3 (Circle packing theorem) Every connected planar graph is the nerve of a circle packing.

Of course, there can be multiple circle packings associated to a given connected planar graph; indeed, since reflections across a line and Möbius transformations map circles to circles (or lines), they will map circle packings to circle packings (unless one or more of the circles is sent to a line). It turns out that once one adds enough edges to the planar graph, the circle packing is otherwise rigid:

Theorem 4 (Koebe-Andreev-Thurston theorem) If a connected planar graph is maximal (i.e., no further edge can be added to it without destroying planarity), then the circle packing given by the above theorem is unique up to reflections and Möbius transformations.

Exercise 5 Let ${G}$ be a connected planar graph with ${n \geq 3}$ vertices. Show that the following are equivalent:

(i) ${G}$ is a maximal planar graph.

(ii) ${G}$ has ${3n-6}$ edges.

(iii) Every drawing ${D}$ of ${G}$ divides the plane into faces that have three edges each. (This includes one unbounded face.)

(iv) At least one drawing ${D}$ of ${G}$ divides the plane into faces that have three edges each.

(Hint: use Euler’s formula ${V-E+F=2}$ , where ${F}$ is the number of faces including the unbounded face.)

Thurston conjectured that circle packings can be used to approximate the conformal map arising in the Riemann mapping theorem. Here is an informal statement:

Conjecture 6 (Informal Thurston conjecture) Let ${U}$ be a simply connected domain, with two distinct points ${z_0,z_1}$ . Let ${\phi: U \rightarrow D(0,1)}$ be the conformal map from ${U}$ to ${D(0,1)}$ that maps ${z_0}$ to the origin and ${z_1}$ to a positive real. For any small ${\varepsilon>0}$ , let ${{\mathcal C}_\varepsilon}$ be the portion of the regular hexagonal circle packing by circles of radius ${\varepsilon}$ that are contained in ${U}$ , and let ${{\mathcal C}'_\varepsilon}$ be an circle packing of ${D(0,1)}$ with all “boundary circles” tangent to ${D(0,1)}$ , giving rise to an “approximate map” ${\phi_\varepsilon: U_\varepsilon \rightarrow D(0,1)}$ defined on the subset ${U_\varepsilon}$ of ${U}$ consisting of the circles of ${{\mathcal C}_\varepsilon}$ , their interiors, and the interstitial regions between triples of mutually tangent circles. Normalise this map so that ${\phi_\varepsilon(z_0)}$ is zero and ${\phi_\varepsilon(z_1)}$ is a positive real. Then ${\phi_\varepsilon}$ converges to ${\phi}$ as ${\varepsilon \rightarrow 0}$ .

A rigorous version of this conjecture was proven by Rodin and Sullivan. Besides some elementary geometric lemmas (regarding the relative sizes of various configurations of tangent circles), the main ingredients are a rigidity result for the regular hexagonal circle packing, and the theory of quasiconformal maps. Quasiconformal maps are what seem on the surface to be a very broad generalisation of the notion of a conformal map. Informally, conformal maps take infinitesimal circles to infinitesimal circles, whereas quasiconformal maps take infinitesimal circles to infinitesimal ellipses of bounded eccentricity. In terms of Wirtinger derivatives, conformal maps obey the Cauchy-Riemann equation ${\frac{\partial \phi}{\partial \overline{z}} = 0}$ , while (sufficiently smooth) quasiconformal maps only obey an inequality ${|\frac{\partial \phi}{\partial \overline{z}}| \leq \frac{K-1}{K+1} |\frac{\partial \phi}{\partial z}|}$ . As such, quasiconformal maps are considerably more plentiful than conformal maps, and in particular it is possible to create piecewise smooth quasiconformal maps by gluing together various simple maps such as affine maps or Möbius transformations; such piecewise maps will naturally arise when trying to rigorously build the map ${\phi_\varepsilon}$ alluded to in the above conjecture. On the other hand, it turns out that quasiconformal maps still have many vestiges of the rigidity properties enjoyed by conformal maps; for instance, there are quasiconformal analogues of fundamental theorems in conformal mapping such as the Schwarz reflection principle, Liouville’s theorem, or Hurwitz’s theorem. Among other things, these quasiconformal rigidity theorems allow one to create conformal maps from the limit of quasiconformal maps in many circumstances, and this will be how the Thurston conjecture will be proven. A key technical tool in establishing these sorts of rigidity theorems will be the theory of an important quasiconformal (quasi-)invariant, the conformal modulus (or, equivalently, the extremal length, which is the reciprocal of the modulus).

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246C notes 1: Meromorphic functions on Riemann surfaces, and the Riemann-Roch theorem

28 March, 2018 in 246C - complex analysis, math.AG, math.CV | Tags: meromorphic function, Riemann surfaces, Riemann-Roch theorem | by Terence Tao | 52 comments

The fundamental object of study in real differential geometry are the real manifolds: Hausdorff topological spaces ${M = M^n}$ that locally look like open subsets of a Euclidean space ${{\bf R}^n}$ , and which can be equipped with an atlas ${(\phi_\alpha: U_\alpha \rightarrow V_\alpha)_{\alpha \in A}}$ of coordinate charts ${\phi_\alpha: U_\alpha \rightarrow V_\alpha}$ from open subsets ${U_\alpha}$ covering ${M}$ to open subsets ${V_\alpha}$ in ${{\bf R}^n}$ , which are homeomorphisms; in particular, the transition maps ${\tau_{\alpha,\beta}: \phi_\alpha( U_\alpha \cap U_\beta ) \rightarrow \phi_\beta( U_\alpha \cap U_\beta )}$ defined by ${\tau_{\alpha,\beta}: \phi_\beta \circ \phi_\alpha^{-1}}$ are all continuous. (It is also common to impose the requirement that the manifold ${M}$ be second countable, though this will not be important for the current discussion.) A smooth real manifold is a real manifold in which the transition maps are all smooth.

In a similar fashion, the fundamental object of study in complex differential geometry are the complex manifolds, in which the model space is ${{\bf C}^n}$ rather than ${{\bf R}^n}$ , and the transition maps ${\tau_{\alpha\beta}}$ are required to be holomorphic (and not merely smooth or continuous). In the real case, the one-dimensional manifolds (curves) are quite simple to understand, particularly if one requires the manifold to be connected; for instance, all compact connected one-dimensional real manifolds are homeomorphic to the unit circle (why?). However, in the complex case, the connected one-dimensional manifolds – the ones that look locally like subsets of ${{\bf C}}$ – are much richer, and are known as Riemann surfaces. For sake of completeness we give the (somewhat lengthy) formal definition:

Definition 1 (Riemann surface) If ${M}$ is a Hausdorff connected topological space, a (one-dimensional complex) atlas is a collection ${(\phi_\alpha: U_\alpha \rightarrow V_\alpha)_{\alpha \in A}}$ of homeomorphisms from open subsets ${(U_\alpha)_{\alpha \in A}}$ of ${M}$ that cover ${M}$ to open subsets ${V_\alpha}$ of the complex numbers ${{\bf C}}$ , such that the transition maps ${\tau_{\alpha,\beta}: \phi_\alpha( U_\alpha \cap U_\beta ) \rightarrow \phi_\beta( U_\alpha \cap U_\beta )}$ defined by ${\tau_{\alpha,\beta}: \phi_\beta \circ \phi_\alpha^{-1}}$ are all holomorphic. Here ${A}$ is an arbitrary index set. Two atlases ${(\phi_\alpha: U_\alpha \rightarrow V_\alpha)_{\alpha \in A}}$ , ${(\phi'_\beta: U'_\beta \rightarrow V'_\beta)_{\beta \in B}}$ on ${M}$ are said to be equivalent if their union is also an atlas, thus the transition maps ${\phi'_\beta \circ \phi_\alpha^{-1}: \phi_\alpha(U_\alpha \cap U'_\beta) \rightarrow \phi'_\beta(U_\alpha \cap U'_\beta)}$ and their inverses are all holomorphic. A Riemann surface is a Hausdorff connected topological space ${M}$ equipped with an equivalence class of one-dimensional complex atlases.

A map ${f: M \rightarrow M'}$ from one Riemann surface ${M}$ to another ${M'}$ is holomorphic if the maps ${\phi'_\beta \circ f \circ \phi_\alpha^{-1}: \phi_\alpha(U_\alpha \cap f^{-1}(U'_\beta)) \rightarrow {\bf C}}$ are holomorphic for any charts ${\phi_\alpha: U_\alpha \rightarrow V_\alpha}$ , ${\phi'_\beta: U'_\beta \rightarrow V'_\beta}$ of an atlas of ${M}$ and ${M'}$ respectively; it is not hard to see that this definition does not depend on the choice of atlas. It is also clear that the composition of two holomorphic maps is holomorphic (and in fact the class of Riemann surfaces with their holomorphic maps forms a category).

Here are some basic examples of Riemann surfaces.

Example 2 (Quotients of ${{\bf C}}$ ) The complex numbers ${{\bf C}}$ clearly form a Riemann surface (using the identity map ${\phi: {\bf C} \rightarrow {\bf C}}$ as the single chart for an atlas). Of course, maps ${f: {\bf C} \rightarrow {\bf C}}$ that are holomorphic in the usual sense will also be holomorphic in the sense of the above definition, and vice versa, so the notion of holomorphicity for Riemann surfaces is compatible with that of holomorphicity for complex maps. More generally, given any discrete additive subgroup ${\Lambda}$ of ${{\bf C}}$ , the quotient ${{\bf C}/\Lambda}$ is a Riemann surface. There are an infinite number of possible atlases to use here; one such is to pick a sufficiently small neighbourhood ${U}$ of the origin in ${{\bf C}}$ and take the atlas ${(\phi_\alpha: U_\alpha \rightarrow U)_{\alpha \in {\bf C}/\Lambda}}$ where ${U_\alpha := \alpha+U}$ and ${\phi_\alpha(\alpha+z) := z}$ for all ${z \in U}$ . In particular, given any non-real complex number ${\omega}$ , the complex torus ${{\bf C} / \langle 1, \omega \rangle}$ formed by quotienting ${{\bf C}}$ by the lattice ${\langle 1, \omega \rangle := \{ n + m \omega: n,m \in {\bf Z}\}}$ is a Riemann surface.

Example 3 Any open connected subset ${U}$ of ${{\bf C}}$ is a Riemann surface. By the Riemann mapping theorem, all simply connected open ${U \subset {\bf C}}$ , other than ${{\bf C}}$ itself, are isomorphic (as Riemann surfaces) to the unit disk (or, equivalently, to the upper half-plane).

Example 4 (Riemann sphere) The Riemann sphere ${{\bf C} \cup \{\infty\}}$ , as a topological manifold, is the one-point compactification of ${{\bf C}}$ . Topologically, this is a sphere and is in particular connected. One can cover the Riemann sphere by the two open sets ${U_1 := {\bf C}}$ and ${U_2 := {\bf C} \cup \{\infty\} \backslash \{0\}}$ , and give these two open sets the charts ${\phi_1: U_1 \rightarrow {\bf C}}$ and ${\phi_2: U_2 \rightarrow {\bf C}}$ defined by ${\phi_1(z) := z}$ for ${z \in {\bf C}}$ , ${\phi_2(z) := 1/z}$ for ${z \in {\bf C} \backslash \{0\}}$ , and ${\phi_2(\infty) := 0}$ . This is a complex atlas since the ${1/z}$ is holomorphic on ${{\bf C} \backslash \{0\}}$ .

An alternate way of viewing the Riemann sphere is as the projective line ${\mathbf{CP}^1}$ . Topologically, this is the punctured complex plane ${{\bf C}^2 \backslash \{(0,0)\}}$ quotiented out by non-zero complex dilations, thus elements of this space are equivalence classes ${[z,w] := \{ (\lambda z, \lambda w): \lambda \in {\bf C} \backslash \{0\}\}}$ with the usual quotient topology. One can cover this space by two open sets ${U_1 := \{ [z,1]: z \in {\bf C} \}}$ and ${U_2: \{ [1,w]: w \in {\bf C} \}}$ and give these two open sets the charts ${\phi: U_1 \rightarrow {\bf C}}$ and ${\phi_2: U_2 \rightarrow {\bf C}}$ defined by ${\phi_1([z,1]) := z}$ for ${z \in {\bf C}}$ , ${\phi_2([1,w]) := w}$ . This is a complex atlas, basically because ${[z,1] = [1,1/z]}$ for ${z \in {\bf C} \backslash \{0\}}$ and ${1/z}$ is holomorphic on ${{\bf C} \backslash \{0\}}$ .

Exercise 5 Verify that the Riemann sphere is isomorphic (as a Riemann surface) to the projective line.

Example 6 (Smooth algebraic plane curves) Let ${P(z_1,z_2,z_3)}$ be a complex polynomial in three variables which is homogeneous of some degree ${d \geq 1}$ , thus

$\displaystyle P( \lambda z_1, \lambda z_2, \lambda z_3) = \lambda^d P( z_1, z_2, z_3). \ \ \ \ \ (1)$

Define the complex projective plane ${\mathbf{CP}^2}$ to be the punctured space ${{\bf C}^3 \backslash \{0\}}$ quotiented out by non-zero complex dilations, with the usual quotient topology. (There is another important topology to place here of fundamental importance in algebraic geometry, namely the Zariski topology, but we will ignore this topology here.) This is a compact space, whose elements are equivalence classes ${[z_1,z_2,z_3] := \{ (\lambda z_1, \lambda z_2, \lambda z_3)\}}$ . Inside this plane we can define the (projective, degree ${d}$ ) algebraic curve

$\displaystyle Z(P) := \{ [z_1,z_2,z_3] \in \mathbf{CP}^2: P(z_1,z_2,z_3) = 0 \};$

this is well defined thanks to (1). It is easy to verify that ${Z(P)}$ is a closed subset of ${\mathbf{CP}^2}$ and hence compact; it is non-empty thanks to the fundamental theorem of algebra.

Suppose that ${P}$ is irreducible, which means that it is not the product of polynomials of smaller degree. As we shall show in the appendix, this makes the algebraic curve connected. (Actually, algebraic curves remain connected even in the reducible case, thanks to Bezout’s theorem, but we will not prove that theorem here.) We will in fact make the stronger nonsingularity hypothesis: there is no triple ${(z_1,z_2,z_3) \in {\bf C}^3 \backslash \{(0,0,0)\}}$ such that the four numbers ${P(z_1,z_2,z_3), \frac{\partial}{\partial z_j} P(z_1,z_2,z_3)}$ simultaneously vanish for ${j=1,2,3}$ . (This looks like four constraints, but is in fact essentially just three, due to the Euler identity

$\displaystyle \sum_{j=1}^3 z_j \frac{\partial}{\partial z_j} P(z_1,z_2,z_3) = d P(z_1,z_2,z_3)$

that arises from differentiating (1) in ${\lambda}$ . The fact that nonsingularity implies irreducibility is another consequence of Bezout’s theorem, which is not proven here.) For instance, the polynomial ${z_1^2 z_3 - z_2^3}$ is irreducible but singular (there is a “cusp” singularity at ${[0,0,1]}$ ). With this hypothesis, we call the curve ${Z(P)}$ smooth.

Now suppose ${[z_1,z_2,z_3]}$ is a point in ${Z(P)}$ ; without loss of generality we may take ${z_3}$ non-zero, and then we can normalise ${z_3=1}$ . Now one can think of ${P(z_1,z_2,1)}$ as an inhomogeneous polynomial in just two variables ${z_1,z_2}$ , and by nondegeneracy we see that the gradient ${(\frac{\partial}{\partial z_1} P(z_1,z_2,1), \frac{\partial}{\partial z_2} P(z_1,z_2,1))}$ is non-zero whenever ${P(z_1,z_2,1)=0}$ . By the (complexified) implicit function theorem, this ensures that the affine algebraic curve

$\displaystyle Z(P)_{aff} := \{ (z_1,z_2) \in {\bf C}^2: P(z_1,z_2,1) = 0 \}$

is a Riemann surface in a neighbourhood of ${(z_1,z_2,1)}$ ; we leave this as an exercise. This can be used to give a coordinate chart for ${Z(P)}$ in a neighbourhood of ${[z_1,z_2,z_3]}$ when ${z_3 \neq 0}$ . Similarly when ${z_1,z_2}$ is non-zero. This can be shown to give an atlas on ${Z(P)}$ , which (assuming the connectedness claim that we will prove later) gives ${Z(P)}$ the structure of a Riemann surface.

Exercise 7 State and prove a complex version of the implicit function theorem that justifies the above claim that the charts in the above example form an atlas, and an algebraic curve associated to a non-singular polynomial is a Riemann surface.

Exercise 8

(i) Show that all (irreducible plane projective) algebraic curves of degree ${1}$ are isomorphic to the Riemann sphere. (Hint: reduce to an explicit linear polynomial such as ${z_3}$ .)

(ii) Show that all (irreducible plane projective) algebraic curves of degree ${2}$ are isomorphic to the Riemann sphere. (Hint: to reduce computation, first use some linear algebra to reduce the homogeneous quadratic polynomial to a standard form, such as ${z_1^2+z_2^2+z_3^2}$ or ${z_2 z_3 - z_1^2}$ .)

Exercise 9 If ${a,b}$ are complex numbers, show that the projective cubic curve

$\displaystyle \{ [z_1, z_2, z_3]: z_2^2 z_3 = z_1^3 + a z_1 z_3^2 + b z_3^3 \}$

is nonsingular if and only if the discriminant ${-16 (4a^3 + 27b^2)}$ is non-zero. (When this occurs, the curve is called an elliptic curve (in Weierstrass form), which is a fundamentally important example of a Riemann surface in many areas of mathematics, and number theory in particular. One can also define the discriminant for polynomials of higher degree, but we will not do so here.)

A recurring theme in mathematics is that an object ${X}$ is often best studied by understanding spaces of “good” functions on ${X}$ . In complex analysis, there are two basic types of good functions:

Definition 10 Let ${X}$ be a Riemann surface. A holomorphic function on ${X}$ is a holomorphic map from ${X}$ to ${{\bf C}}$ ; the space of all such functions will be denoted ${{\mathcal O}(X)}$ . A meromorphic function on ${X}$ is a holomorphic map from ${X}$ to the Riemann sphere ${{\bf C} \cup \{\infty\}}$ , that is not identically equal to ${\infty}$ ; the space of all such functions will be denoted ${M(X)}$ .

One can also define holomorphicity and meromorphicity in terms of charts: a function ${f: X \rightarrow {\bf C}}$ is holomorphic if and only if, for any chart ${\phi_\alpha: U_\alpha \rightarrow {\bf C}}$ , the map ${f \circ \phi^{-1}_\alpha: \phi_\alpha(U_\alpha) \rightarrow {\bf C}}$ is holomorphic in the usual complex analysis sense; similarly, a function ${f: X \rightarrow {\bf C} \cup \{\infty\}}$ is meromorphic if and only if the preimage ${f^{-1}(\{\infty\})}$ is discrete (otherwise, by analytic continuation and the connectedness of ${X}$ , ${f}$ will be identically equal to ${\infty}$ ) and for any chart ${\phi_\alpha: U_\alpha \rightarrow X}$ , the map ${f \circ \phi_\alpha^{-1}: \phi_\alpha(U_\alpha) \rightarrow {\bf C} \cup \{\infty\}}$ becomes a meromorphic function in the usual complex analysis sense, after removing the discrete set of complex numbers where this map is infinite. One consequence of this alternate definition is that the space ${{\mathcal O}(X)}$ of holomorphic functions is a commutative complex algebra (a complex vector space closed under pointwise multiplication), while the space ${M(X)}$ of meromorphic functions is a complex field (a commutative complex algebra where every non-zero element has an inverse). Another consequence is that one can define the notion of a zero of given order ${k}$ , or a pole of order ${k}$ , for a holomorphic or meromorphic function, by composing with a chart map and using the usual complex analysis notions there, noting (from the holomorphicity of transition maps and their inverses) that this does not depend on the choice of chart. (However, one cannot similarly define the residue of a meromorphic function on ${X}$ this way, as the residue turns out to be chart-dependent thanks to the chain rule. Residues should instead be applied to meromorphic ${1}$ -forms, a concept we will introduce later.) A third consequence is analytic continuation: if two holomorphic or meromorphic functions on ${X}$ agree on a non-empty open set, then they agree everywhere.

On the complex numbers ${{\bf C}}$ , there are of course many holomorphic functions and meromorphic functions; for instance any power series with an infinite radius of convergence will give a holomorphic function, and the quotient of any two such functions (with non-zero denominator) will give a meromorphic function. Furthermore, we have extremely wide latitude in how to specify the zeroes of the holomorphic function, or the zeroes and poles of the meromorphic function, thanks to tools such as the Weierstrass factorisation theorem or the Mittag-Leffler theorem (covered in previous quarters).

It turns out, however, that the situation changes dramatically when the Riemann surface ${X}$ is compact, with the holomorphic and meromorphic functions becoming much more rigid. First of all, compactness eliminates all holomorphic functions except for the constants:

Lemma 11 Let ${f \in \mathcal{O}(X)}$ be a holomorphic function on a compact Riemann surface ${X}$ . Then ${f}$ is constant.

This result should be seen as a close sibling of Liouville’s theorem that all bounded entire functions are constant. (Indeed, in the case of a complex torus, this lemma is a corollary of Liouville’s theorem.)

Proof: As ${f}$ is continuous and ${X}$ is compact, ${|f(z_0)|}$ must attain a maximum at some point ${z_0 \in X}$ . Working in a chart around ${z_0}$ and applying the maximum principle, we conclude that ${f}$ is constant in a neighbourhood of ${z_0}$ , and hence is constant everywhere by analytic continuation. $\Box$

This dramatically cuts down the number of possible meromorphic functions – indeed, for an abstract Riemann surface, it is not immediately obvious that there are any non-constant meromorphic functions at all! As the poles are isolated and the surface is compact, a meromorphic function can only have finitely many poles, and if one prescribes the location of the poles and the maximum order at each pole, then we shall see that the space of meromorphic functions is now finite dimensional. The precise dimensions of these spaces are in fact rather interesting, and obey a basic duality law known as the Riemann-Roch theorem. We will give a mostly self-contained proof of the Riemann-Roch theorem in these notes, omitting only some facts about genus and Euler characteristic, as well as construction of certain meromorphic ${1}$ -forms (also known as Abelian differentials).

A more detailed study of Riemann surface (and more generally, complex manifolds) can be found for instance in Griffiths and Harris’s “Principles of Algebraic Geometry“.

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Course announcement: 246C, complex analysis

13 March, 2018 in 246C - complex analysis, advertising, math.CV | by Terence Tao | 8 comments

Next quarter (starting Monday, April 2) I will be teaching Math 246C (complex analysis) here at UCLA. This is the third in a three-series graduate course on complex analysis; a few years ago I taught the first course in this series (246A), so this course can be thought of in some sense as a sequel to that one (and would certainly assume knowledge of the material in that course as a prerequisite), although it also assumes knowledge of material from the second course 246B (which covers such topics as Weierstrass factorization and the theory of harmonic functions).

246C is primarily a topics course, and tends to be a somewhat miscellaneous collection of complex analysis subjects that were not covered in the previous two installments of the series. The initial topics I have in mind to cover are

Elliptic functions;
The Riemann-Roch theorem;
Circle packings;
The Bieberbach conjecture (proven by de Branges); and
the Schramm-Loewner equation (SLE).
This list is however subject to change (it is the first time I will have taught on any of these topics, and I am not yet certain on the most logical way to arrange them; also I am not completely certain that I will be able to cover all the above topics in ten weeks). I welcome reference recommendations and other suggestions from readers who have taught on one or more of these topics.

As usual, I will be posting lecture notes on this blog as the course progresses.

[Update: Mar 13: removed elliptic functions, as I have just learned that this was already covered in the prior 246B course.]

Heat flow and zeroes of polynomials

17 October, 2017 in expository, math.AP, math.CV, math.DS | Tags: entropy, heat flow, polynomials, zeroes | by Terence Tao | 34 comments

Let ${P(z) = z^n + a_{n-1} z^{n-1} + \dots + a_0}$ be a monic polynomial of degree ${n}$ with complex coefficients. Then by the fundamental theorem of algebra, we can factor ${P}$ as

$\displaystyle P(z) = (z-z_1) \dots (z-z_n) \ \ \ \ \ (1)$

for some complex zeroes ${z_1,\dots,z_n}$ (possibly with repetition).

Now suppose we evolve ${P}$ with respect to time by heat flow, creating a function ${P(t,z)}$ of two variables with given initial data ${P(0,z) = P(z)}$ for which

$\displaystyle \partial_t P(t,z) = \partial_{zz} P(t,z). \ \ \ \ \ (2)$

On the space of polynomials of degree at most ${n}$ , the operator ${\partial_{zz}}$ is nilpotent, and one can solve this equation explicitly both forwards and backwards in time by the Taylor series

$\displaystyle P(t,z) = \sum_{j=0}^\infty \frac{t^j}{j!} \partial_z^{2j} P(0,z).$

For instance, if one starts with a quadratic ${P(0,z) = z^2 + bz + c}$ , then the polynomial evolves by the formula

$\displaystyle P(t,z) = z^2 + bz + (c+2t).$

As the polynomial ${P(t)}$ evolves in time, the zeroes ${z_1(t),\dots,z_n(t)}$ evolve also. Assuming for sake of discussion that the zeroes are simple, the inverse function theorem tells us that the zeroes will (locally, at least) evolve smoothly in time. What are the dynamics of this evolution?

For instance, in the quadratic case, the quadratic formula tells us that the zeroes are

$\displaystyle z_1(t) = \frac{-b + \sqrt{b^2 - 4(c+2t)}}{2}$

and

$\displaystyle z_2(t) = \frac{-b - \sqrt{b^2 - 4(c+2t)}}{2}$

after arbitrarily choosing a branch of the square root. If ${b,c}$ are real and the discriminant ${b^2 - 4c}$ is initially positive, we see that we start with two real zeroes centred around ${-b/2}$ , which then approach each other until time ${t = \frac{b^2-4c}{8}}$ , at which point the roots collide and then move off from each other in an imaginary direction.

In the general case, we can obtain the equations of motion by implicitly differentiating the defining equation

$\displaystyle P( t, z_i(t) ) = 0$

in time using (2) to obtain

$\displaystyle \partial_{zz} P( t, z_i(t) ) + \partial_t z_i(t) \partial_z P(t,z_i(t)) = 0.$

To simplify notation we drop the explicit dependence on time, thus

$\displaystyle \partial_{zz} P(z_i) + (\partial_t z_i) \partial_z P(z_i)= 0.$

From (1) and the product rule, we see that

$\displaystyle \partial_z P( z_i ) = \prod_{j:j \neq i} (z_i - z_j)$

and

$\displaystyle \partial_{zz} P( z_i ) = 2 \sum_{k:k \neq i} \prod_{j:j \neq i,k} (z_i - z_j)$

(where all indices are understood to range over ${1,\dots,n}$ ) leading to the equations of motion

$\displaystyle \partial_t z_i = \sum_{k:k \neq i} \frac{2}{z_k - z_i}, \ \ \ \ \ (3)$

at least when one avoids those times in which there is a repeated zero. In the case when the zeroes ${z_i}$ are real, each term ${\frac{2}{z_k-z_i}}$ represents a (first-order) attraction in the dynamics between ${z_i}$ and ${z_k}$ , but the dynamics are more complicated for complex zeroes (e.g. purely imaginary zeroes will experience repulsion rather than attraction, as one already sees in the quadratic example). Curiously, this system resembles that of Dyson brownian motion (except with the brownian motion part removed, and time reversed). I learned of the connection between the ODE (3) and the heat equation from this paper of Csordas, Smith, and Varga, but perhaps it has been mentioned in earlier literature as well.

One interesting consequence of these equations is that if the zeroes are real at some time, then they will stay real as long as the zeroes do not collide. Let us now restrict attention to the case of real simple zeroes, in which case we will rename the zeroes as ${x_i}$ instead of ${z_i}$ , and order them as ${x_1 < \dots < x_n}$ . The evolution

$\displaystyle \partial_t x_i = \sum_{k:k \neq i} \frac{2}{x_k - x_i}$

can now be thought of as reverse gradient flow for the “entropy”

$\displaystyle H := -\sum_{i,j: i \neq j} \log |x_i - x_j|,$

(which is also essentially the logarithm of the discriminant of the polynomial) since we have

$\displaystyle \partial_t x_i = \frac{\partial H}{\partial x_i}.$

In particular, we have the monotonicity formula

$\displaystyle \partial_t H = 4E$

where ${E}$ is the “energy”

$\displaystyle E := \frac{1}{4} \sum_i (\frac{\partial H}{\partial x_i})^2$

$\displaystyle = \sum_i (\sum_{k:k \neq i} \frac{1}{x_k-x_i})^2$

$\displaystyle = \sum_{i,k: i \neq k} \frac{1}{(x_k-x_i)^2} + 2 \sum_{i,j,k: i,j,k \hbox{ distinct}} \frac{1}{(x_k-x_i)(x_j-x_i)}$

$\displaystyle = \sum_{i,k: i \neq k} \frac{1}{(x_k-x_i)^2}$

where in the last line we use the antisymmetrisation identity

$\displaystyle \frac{1}{(x_k-x_i)(x_j-x_i)} + \frac{1}{(x_i-x_j)(x_k-x_j)} + \frac{1}{(x_j-x_k)(x_i-x_k)} = 0.$

Among other things, this shows that as one goes backwards in time, the entropy decreases, and so no collisions can occur to the past, only in the future, which is of course consistent with the attractive nature of the dynamics. As ${H}$ is a convex function of the positions ${x_1,\dots,x_n}$ , one expects ${H}$ to also evolve in a convex manner in time, that is to say the energy ${E}$ should be increasing. This is indeed the case:

Exercise 1 Show that

$\displaystyle \partial_t E = 2 \sum_{i,j: i \neq j} (\frac{2}{(x_i-x_j)^2} - \sum_{k: i,j,k \hbox{ distinct}} \frac{1}{(x_k-x_i)(x_k-x_j)})^2.$

Symmetric polynomials of the zeroes are polynomial functions of the coefficients and should thus evolve in a polynomial fashion. One can compute this explicitly in simple cases. For instance, the center of mass is an invariant:

$\displaystyle \partial_t \frac{1}{n} \sum_i x_i = 0.$

The variance decreases linearly:

Exercise 2 Establish the virial identity

$\displaystyle \partial_t \sum_{i,j} (x_i-x_j)^2 = - 4n^2(n-1).$

As the variance (which is proportional to ${\sum_{i,j} (x_i-x_j)^2}$ ) cannot become negative, this identity shows that “finite time blowup” must occur – that the zeroes must collide at or before the time ${\frac{1}{4n^2(n-1)} \sum_{i,j} (x_i-x_j)^2}$ .

Exercise 3 Show that the Stieltjes transform

$\displaystyle s(t,z) = \sum_i \frac{1}{x_i - z}$

solves the viscous Burgers equation

$\displaystyle \partial_t s = \partial_{zz} s - 2 s \partial_z s,$

either by using the original heat equation (2) and the identity ${s = - \partial_z P / P}$ , or else by using the equations of motion (3). This relation between the Burgers equation and the heat equation is known as the Cole-Hopf transformation.

The paper of Csordas, Smith, and Varga mentioned previously gives some other bounds on the lifespan of the dynamics; roughly speaking, they show that if there is one pair of zeroes that are much closer to each other than to the other zeroes then they must collide in a short amount of time (unless there is a collision occuring even earlier at some other location). Their argument extends also to situations where there are an infinite number of zeroes, which they apply to get new results on Newman’s conjecture in analytic number theory. I would be curious to know of further places in the literature where this dynamics has been studied.

Another problem about power series

23 October, 2016 in 246A - complex analysis, math.CV, math.RA, question | Tags: identity, power series | by Terence Tao | 37 comments

By an odd coincidence, I stumbled upon a second question in as many weeks about power series, and once again the only way I know how to prove the result is by complex methods; once again, I am leaving it here as a challenge to any interested readers, and I would be particularly interested in knowing of a proof that was not based on complex analysis (or thinly disguised versions thereof), or for a reference to previous literature where something like this identity has occured. (I suspect for instance that something like this may have shown up before in free probability, based on the answer to part (ii) of the problem.)

Here is a purely algebraic form of the problem:

Problem 1 Let ${F = F(z)}$ be a formal function of one variable ${z}$ . Suppose that ${G = G(z)}$ is the formal function defined by

$\displaystyle G := \sum_{n=1}^\infty \left( \frac{F^n}{n!} \right)^{(n-1)}$

$\displaystyle = F + \left(\frac{F^2}{2}\right)' + \left(\frac{F^3}{6}\right)'' + \dots$

$\displaystyle = F + FF' + (F (F')^2 + \frac{1}{2} F^2 F'') + \dots,$

where we use ${f^{(k)}}$ to denote the ${k}$ -fold derivative of ${f}$ with respect to the variable ${z}$ .

(i) Show that ${F}$ can be formally recovered from ${G}$ by the formula
$\displaystyle F = \sum_{n=1}^\infty (-1)^{n-1} \left( \frac{G^n}{n!} \right)^{(n-1)}$

$\displaystyle = G - \left(\frac{G^2}{2}\right)' + \left(\frac{G^3}{6}\right)'' - \dots$

$\displaystyle = G - GG' + (G (G')^2 + \frac{1}{2} G^2 G'') - \dots.$

(ii) There is a remarkable further formal identity relating ${F(z)}$ with ${G(z)}$ that does not explicitly involve any infinite summation. What is this identity?

To rigorously formulate part (i) of this problem, one could work in the commutative differential ring of formal infinite series generated by polynomial combinations of ${F}$ and its derivatives (with no constant term). Part (ii) is a bit trickier to formulate in this abstract ring; the identity in question is easier to state if ${F, G}$ are formal power series, or (even better) convergent power series, as it involves operations such as composition or inversion that can be more easily defined in those latter settings.

To illustrate Problem 1(i), let us compute up to third order in ${F}$ , using ${{\mathcal O}(F^4)}$ to denote any quantity involving four or more factors of ${F}$ and its derivatives, and similarly for other exponents than ${4}$ . Then we have

$\displaystyle G = F + FF' + (F (F')^2 + \frac{1}{2} F^2 F'') + {\mathcal O}(F^4)$

and hence

$\displaystyle G' = F' + (F')^2 + FF'' + {\mathcal O}(F^3)$

$\displaystyle G'' = F'' + {\mathcal O}(F^2);$

multiplying, we have

$\displaystyle GG' = FF' + F (F')^2 + F^2 F'' + F (F')^2 + {\mathcal O}(F^4)$

and

$\displaystyle G (G')^2 + \frac{1}{2} G^2 G'' = F (F')^2 + \frac{1}{2} F^2 F'' + {\mathcal O}(F^4)$

and hence after a lot of canceling

$\displaystyle G - GG' + (G (G')^2 + \frac{1}{2} G^2 G'') = F + {\mathcal O}(F^4).$

Thus Problem 1(i) holds up to errors of ${{\mathcal O}(F^4)}$ at least. In principle one can continue verifying Problem 1(i) to increasingly high order in ${F}$ , but the computations rapidly become quite lengthy, and I do not know of a direct way to ensure that one always obtains the required cancellation at the end of the computation.

Problem 1(i) can also be posed in formal power series: if

$\displaystyle F(z) = a_1 z + a_2 z^2 + a_3 z^3 + \dots$

is a formal power series with no constant term with complex coefficients ${a_1, a_2, \dots}$ with ${|a_1|<1}$ , then one can verify that the series

$\displaystyle G := \sum_{n=1}^\infty \left( \frac{F^n}{n!} \right)^{(n-1)}$

makes sense as a formal power series with no constant term, thus

$\displaystyle G(z) = b_1 z + b_2 z^2 + b_3 z^3 + \dots.$

For instance it is not difficult to show that ${b_1 = \frac{a_1}{1-a_1}}$ . If one further has ${|b_1| < 1}$ , then it turns out that

$\displaystyle F = \sum_{n=1}^\infty (-1)^{n-1} \left( \frac{G^n}{n!} \right)^{(n-1)}$

as formal power series. Currently the only way I know how to show this is by first proving the claim for power series with a positive radius of convergence using the Cauchy integral formula, but even this is a bit tricky unless one has managed to guess the identity in (ii) first. (In fact, the way I discovered this problem was by first trying to solve (a variant of) the identity in (ii) by Taylor expansion in the course of attacking another problem, and obtaining the transform in Problem 1 as a consequence.)

The transform that takes ${F}$ to ${G}$ resembles both the exponential function

$\displaystyle \exp(F) = \sum_{n=0}^\infty \frac{F^n}{n!}$

and Taylor’s formula

$\displaystyle F(z) = \sum_{n=0}^\infty \frac{F^{(n)}(0)}{n!} z^n$

but does not seem to be directly connected to either (this is more apparent once one knows the identity in (ii)).

246A, Notes 5: conformal mapping

18 October, 2016 in 246A - complex analysis, math.CV, math.DG | Tags: conformal mapping, Green's function, Riemann mapping theorem, Schwarz-Christoffel formula, uniformisation theorem | by Terence Tao | 31 comments

In the previous set of notes we introduced the notion of a complex diffeomorphism ${f: U \rightarrow V}$ between two open subsets ${U,V}$ of the complex plane ${{\bf C}}$ (or more generally, two Riemann surfaces): an invertible holomorphic map whose inverse was also holomorphic. (Actually, the last part is automatic, thanks to Exercise 40 of Notes 4.) Such maps are also known as biholomorphic maps or conformal maps (although in some literature the notion of “conformal map” is expanded to permit maps such as the complex conjugation map ${z \mapsto \overline{z}}$ that are angle-preserving but not orientation-preserving, as well as maps such as the exponential map ${z \mapsto \exp(z)}$ from ${{\bf C}}$ to ${{\bf C} \backslash \{0\}}$ that are only locally injective rather than globally injective). Such complex diffeomorphisms can be used in complex analysis (or in the analysis of harmonic functions) to change the underlying domain ${U}$ to a domain that may be more convenient for calculations, thanks to the following basic lemma:

Lemma 1 (Holomorphicity and harmonicity are conformal invariants) Let ${\phi: U \rightarrow V}$ be a complex diffeomorphism between two Riemann surfaces ${U,V}$ .

(i) If ${f: V \rightarrow W}$ is a function to another Riemann surface ${W}$ , then ${f}$ is holomorphic if and only if ${f \circ \phi: U \rightarrow W}$ is holomorphic.

(ii) If ${U,V}$ are open subsets of ${{\bf C}}$ and ${u: V \rightarrow {\bf R}}$ is a function, then ${u}$ is harmonic if and only if ${u \circ \phi: U \rightarrow {\bf R}}$ is harmonic.

Proof: Part (i) is immediate since the composition of two holomorphic functions is holomorphic. For part (ii), observe that if ${u: V \rightarrow {\bf R}}$ is harmonic then on any ball ${B(z_0,r)}$ in ${V}$ , ${u}$ is the real part of some holomorphic function ${f: B(z_0,r) \rightarrow {\bf C}}$ thanks to Exercise 62 of Notes 3. By part (i), ${f \circ \phi: B(z_0,r) \rightarrow {\bf C}}$ is also holomorphic. Taking real parts we see that ${u \circ \phi}$ is harmonic on each ball ${B(z_0,r)}$ in ${V}$ , and hence harmonic on all of ${V}$ , giving one direction of (ii); the other direction is proven similarly. $\Box$

Exercise 2 Establish Lemma 1(ii) by direct calculation, avoiding the use of holomorphic functions. (Hint: the calculations are cleanest if one uses Wirtinger derivatives, as per Exercise 27 of Notes 1.)

Exercise 3 Let ${\phi: U \rightarrow V}$ be a complex diffeomorphism between two open subsets ${U,V}$ of ${{\bf C}}$ , let ${z_0}$ be a point in ${U}$ , let ${m}$ be a natural number, and let ${f: V \rightarrow {\bf C} \cup \{\infty\}}$ be holomorphic. Show that ${f: V \rightarrow {\bf C} \cup \{\infty\}}$ has a zero (resp. a pole) of order ${m}$ at ${\phi(z_0)}$ if and only if ${f \circ \phi: U \rightarrow {\bf C} \cup \{\infty\}}$ has a zero (resp. a pole) of order ${m}$ at ${z_0}$ .

From Lemma 1(ii) we can now define the notion of a harmonic function ${u: M \rightarrow {\bf R}}$ on a Riemann surface ${M}$ ; such a function ${u}$ is harmonic if, for every coordinate chart ${\phi_\alpha: U_\alpha \rightarrow V_\alpha}$ in some atlas, the map ${u \circ \phi_\alpha^{-1}: V_\alpha \rightarrow {\bf R}}$ is harmonic. Lemma 1(ii) ensures that this definition of harmonicity does not depend on the choice of atlas. Similarly, using Exercise 3 one can define what it means for a holomorphic map ${f: M \rightarrow {\bf C} \cup \{\infty\}}$ on a Riemann surface ${M}$ to have a pole or zero of a given order at a point ${p_0 \in M}$ , with the definition being independent of the choice of atlas.

In view of Lemma 1, it is thus natural to ask which Riemann surfaces are complex diffeomorphic to each other, and more generally to understand the space of holomorphic maps from one given Riemann surface to another. We will initially focus attention on three important model Riemann surfaces:

(i) (Elliptic model) The Riemann sphere ${{\bf C} \cup \{\infty\}}$ ;
(ii) (Parabolic model) The complex plane ${{\bf C}}$ ; and
(iii) (Hyperbolic model) The unit disk ${D(0,1)}$ .

The designation of these model Riemann surfaces as elliptic, parabolic, and hyperbolic comes from Riemannian geometry, where it is natural to endow each of these surfaces with a constant curvature Riemannian metric which is positive, zero, or negative in the elliptic, parabolic, and hyperbolic cases respectively. However, we will not discuss Riemannian geometry further here.

All three model Riemann surfaces are simply connected, but none of them are complex diffeomorphic to any other; indeed, there are no non-constant holomorphic maps from the Riemann sphere to the plane or the disk, nor are there any non-constant holomorphic maps from the plane to the disk (although there are plenty of holomorphic maps going in the opposite directions). The complex automorphisms (that is, the complex diffeomorphisms from a surface to itself) of each of the three surfaces can be classified explicitly. The automorphisms of the Riemann sphere turn out to be the Möbius transformations ${z \mapsto \frac{az+b}{cz+d}}$ with ${ad-bc \neq 0}$ , also known as fractional linear transformations. The automorphisms of the complex plane are the linear transformations ${z \mapsto az+b}$ with ${a \neq 0}$ , and the automorphisms of the disk are the fractional linear transformations of the form ${z \mapsto e^{i\theta} \frac{\alpha - z}{1 - \overline{\alpha} z}}$ for ${\theta \in {\bf R}}$ and ${\alpha \in D(0,1)}$ . Holomorphic maps ${f: D(0,1) \rightarrow D(0,1)}$ from the disk ${D(0,1)}$ to itself that fix the origin obey a basic but incredibly important estimate known as the Schwarz lemma: they are “dominated” by the identity function ${z \mapsto z}$ in the sense that ${|f(z)| \leq |z|}$ for all ${z \in D(0,1)}$ . Among other things, this lemma gives guidance to determine when a given Riemann surface is complex diffeomorphic to a disk; we shall discuss this point further below.

It is a beautiful and fundamental fact in complex analysis that these three model Riemann surfaces are in fact an exhaustive list of the simply connected Riemann surfaces, up to complex diffeomorphism. More precisely, we have the Riemann mapping theorem and the uniformisation theorem:

Theorem 4 (Riemann mapping theorem) Let ${U}$ be a simply connected open subset of ${{\bf C}}$ that is not all of ${{\bf C}}$ . Then ${U}$ is complex diffeomorphic to ${D(0,1)}$ .

Theorem 5 (Uniformisation theorem) Let ${M}$ be a simply connected Riemann surface. Then ${M}$ is complex diffeomorphic to ${{\bf C} \cup \{\infty\}}$ , ${{\bf C}}$ , or ${D(0,1)}$ .

As we shall see, every connected Riemann surface can be viewed as the quotient of its simply connected universal cover by a discrete group of automorphisms known as deck transformations. This in principle gives a complete classification of Riemann surfaces up to complex diffeomorphism, although the situation is still somewhat complicated in the hyperbolic case because of the wide variety of discrete groups of automorphisms available in that case.

We will prove the Riemann mapping theorem in these notes, using the elegant argument of Koebe that is based on the Schwarz lemma and Montel’s theorem (Exercise 57 of Notes 4). The uniformisation theorem is however more difficult to establish; we discuss some components of a proof (based on the Perron method of subharmonic functions) here, but stop short of providing a complete proof.

The above theorems show that it is in principle possible to conformally map various domains into model domains such as the unit disk, but the proofs of these theorems do not readily produce explicit conformal maps for this purpose. For some domains we can just write down a suitable such map. For instance:

Exercise 6 (Cayley transform) Let ${{\bf H} := \{ z \in {\bf C}: \mathrm{Im} z > 0 \}}$ be the upper half-plane. Show that the Cayley transform ${\phi: {\bf H} \rightarrow D(0,1)}$ , defined by

$\displaystyle \phi(z) := \frac{z-i}{z+i},$

is a complex diffeomorphism from the upper half-plane ${{\bf H}}$ to the disk ${D(0,1)}$ , with inverse map ${\phi^{-1}: D(0,1) \rightarrow {\bf H}}$ given by

$\displaystyle \phi^{-1}(w) := i \frac{1+w}{1-w}.$

Exercise 7 Show that for any real numbers ${a<b}$ , the strip ${\{ z \in {\bf C}: a < \mathrm{Re}(z) < b \}}$ is complex diffeomorphic to the disk ${D(0,1)}$ . (Hint: use the complex exponential and a linear transformation to map the strip onto the half-plane ${{\bf H}}$ .)

Exercise 8 Show that for any real numbers ${a<b<a+2\pi}$ , the strip ${\{ re^{i\theta}: r>0, a < \theta < b \}}$ is complex diffeomorphic to the disk ${D(0,1)}$ . (Hint: use a branch of either the complex logarithm, or of a complex power ${z \mapsto z^\alpha}$ .)

We will discuss some other explicit conformal maps in this set of notes, such as the Schwarz-Christoffel maps that transform the upper half-plane ${{\bf H}}$ to polygonal regions. Further examples of conformal mapping can be found in the text of Stein-Shakarchi.

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Category Archive

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